CS 296.

1
Mathematical Modelling of Continuous Systems
Carlo Tomasi
Duke University
Fall 2004
2
Chapter 1
Introduction
Fields such as robotics or computer vision are interdisciplinary subjects at the intersection of engineering and computer
science. By their nature, they deal with both computers and the physical world. Although the former are in the latter,
the workings of computers are best described in the black-and-white vocabulary of discrete mathematics, which is
foreign to most classical models of reality, quantum physics notwithstanding.
This class surveys some of the key tools of applied math to be used at the interface of continuous and discrete. It
is not on robotics or computer vision, nor does it cover any other application area. Applications evolve rapidly, but
their mathematical foundations remain. Even if you will not pursue any of these fields, the mathematics that you learn
in this class will not go wasted. To be sure, applied mathematics is a discipline in itself and, in many universities, a
separate department. Consequently, this class can be a quick tour at best. It does not replace calculus or linear algebra,
which are assumed as prerequisites, nor is it a comprehensive survey of applied mathematics. What is covered is a
compromise between the time available and what is useful and fun to talk about. Even if in some cases you may have
to wait until you take an applied class to fully appreciate the usefulness of a particular topic, I hope that you will enjoy
studying these subjects in their own right.
1.1 Who Should Take This Class
The main goal of this class is to present a collection of mathematical tools for both understanding and solving problems
in fields that manipulate models of the real world, such as robotics, artificial intelligence, vision, engineering, or several
aspects of the biological sciences. Several classes at most universities each cover some of the topics presented in this
class, and do so in much greater detail. If you want to understand the full details of any one of the topics in the
syllabus below, you should take one or more of these other classes instead. If you want to understand how these tools
are implemented numerically, you should take one of the classes in the scientific computing program, which again
cover these issues in much better detail. Finally, if you want to understand robotics, vision, or other applied fields, you
should take classes in these subjects, since this course is not on applications.
On the other hand, if you do plan to study robotics, vision, or other applied subjects in the future, and you regard
yourself as a user of the mathematical techniques outlined in the syllabus below, then you may benefit from this course.
Of the proofs, we will only see those that add understanding. Of the implementation aspects of algorithms that are
available in, say, Matlab or LApack, we will only see the parts that we need to understand when we use the code.
In brief, we will be able to cover more topics than other classes because we will be often (but not always) un-
concerned with rigorous proof or implementation issues. The emphasis will be on intuition and on practicality of the
various algorithms. For instance, why are singular values important, and how do they relate to eigenvalues? What are
the dangers of Newton-style minimization? How does a Kalman filter work, and why do PDEs lead to sparse linear
systems? In this spirit, for instance, we discuss Singular Value Decomposition and Schur decomposition both because
they never fail and because they clarify the structure of an algebraic or a differential linear problem.
3
4 CHAPTER 1. INTRODUCTION
1.2 Syllabus
Here is the ideal syllabus, but how much we cover depends on how fast we go.
1. Introduction
2. Unknown numbers
2.1 Algebraic linear systems
2.1.1 Characterization of the solutions to a linear system
2.1.2 Gaussian elimination
2.1.3 The Singular Value Decomposition
2.1.4 The pseudoinverse
2.2 Function optimization
2.2.1 Newton and Gauss-Newton methods
2.2.2 Levenberg-Marquardt method
2.2.3 Constraints and Lagrange multipliers
3. Unknown functions of one real variable
3.1 Ordinary differential linear systems
3.1.1 Eigenvalues and eigenvectors
3.1.2 The Schur decomposition
3.1.3 Ordinary differential linear systems
3.1.4 The matrix zoo
3.1.5 Real, symmetric, positive-definite matrices
3.2 Statistical estimation
3.2.1 Linear estimation
3.2.2 Weighted least squares
3.2.3 The Kalman filter
4. Unknown functions of several variables
4.1 Tensor fields of several variables
4.1.1 Grad, div, curl
4.1.2 Line, surface, and volume integrals
4.1.3 Green’s theorem and potential fields of two variables
4.1.4 Stokes’ and divergence theorems and potential fields of three variables
4.1.5 Diffusion and flow problems
4.2 Partial differential equations and sparse linear systems
4.2.1 Finite differences
4.2.2 Direct versus iterative solution methods
4.2.3 Jacobi and Gauss-Seidel iterations
4.2.4 Successive overrelaxation
4.3 Calculus of variations
4.3.1 Euler-Lagrange equations
4.3.2 The brachistochrone
1.3. DISCUSSION OF THE SYLLABUS 5
1.3 Discussion of the Syllabus
In robotics, vision, physics, and any other branch of science whose subject belongs to or interacts with the real
world, mathematical models are developed that describe the relationship between different quantities. Some of these
quantities are measured, or sensed, while others are inferred by calculation. For instance, in computer vision, equations
tie the coordinates of points in space to the coordinates of corresponding points in different images. Image points are
data, world points are unknowns to be computed.
Similarly, in robotics, a robot arm is modeled by equations that describe where each link of the robot is as a
function of the configuration of the link’s own joints and that of the links that support it. The desired position of the
end effector, as well as the current configuration of all the joints, are the data. The unknowns are the motions to be
imparted to the joints so that the end effector reaches the desired target position.
Of course, what is data and what is unknown depends on the problem. For instance, the vision system mentioned
above could be looking at the robot arm. Then, the robot’s end effector position could be the unknowns to be solved
for by the vision system. Once vision has solved its problem, it could feed the robot’s end-effector position as data for
the robot controller to use in its own motion planning problem.
Sensed data are invariably noisy, because sensors have inherent limitations of accuracy, precision, resolution,
and repeatability. Consequently, the systems of equations to be solved are typically overconstrained: there are more
equations than unknowns, and it is hoped that the errors that affect the coefficients of one equation are partially
cancelled by opposite errors in other equations. This is the basis of optimization problems: Rather than solving a
minimal system exactly, an optimization problem tries to solve many equations simultaneously, each of them only
approximately, but collectively as well as possible, according to some global criterion. Least squares is perhaps the
most popular such criterion, and we will devote a good deal of attention to it.
In summary, the problems encountered in robotics and vision, as well as other applications of mathematics, are
optimization problems. A fundamental distinction between different classes of problems reflects the complexity of the
unknowns. In the simplest case, unknowns are scalars. When there is more than one scalar, the unknown is a vector
of numbers, typically either real or complex. Accordingly, the first part of this course will be devoted to describing
systems of algebraic equations, especially linear equations, and optimization techniques for problems whose solution
is a vector of reals. The main tool for understanding linear algebraic systems is the Singular Value Decomposition
(SVD), which is both conceptually fundamental and practically of extreme usefulness. When the systems are nonlinear,
they can be solved by various techniques of function optimization, of which we will consider the basic aspects.
Since physical quantities often evolve over time, many problems arise in which the unknowns are themselves
functions of time. This is our second class of problems. Again, problems can be cast as a set of equations to be solved
exactly, and this leads to the theory of Ordinary Differential Equations (ODEs). Here, “ordinary” expresses the fact
that the unknown functions depend on just one variable (e.g., time). The main conceptual tool for addressing ODEs is
the theory of eigenvalues, and the primary computational tool is the Schur decomposition.
Alternatively, problems with time varying solutions can be stated as minimization problems. When viewed glob-
ally, these minimization problems lead to the calculus of variations. When the minimization problems above are
studied locally, they become state estimation problems, and the relevant theory is that of dynamic systems and Kalman
filtering.
The third category of problems concerns unknown functions of more than one variable. The images taken by a
moving camera, for instance, are functions of time and space, and so are the unknown quantities that one can compute
fromthe images, such as the distance of points in the world fromthe camera. This leads to Partial Differential equations
(PDEs), or to extensions of the calculus of variations. In this class, we will see how PDEs arise, and how they can be
solved numerically.
6 CHAPTER 1. INTRODUCTION
1.4 Books
The class will be based on these lecture notes, and additional notes handed out when necessary. Other useful references
include the following.
• R. Courant and D. Hilbert, Methods of Mathematical Physics, Volume I and II, John Wiley and Sons, 1989.
• D. A. Danielson, Vectors and Tensors in Engineering and Physics, Addison-Wesley, 1992.
• J. W. Demmel, Applied Numerical Linear Algebra, SIAM, 1997.
• A. Gelb et al., Applied Optimal Estimation, MIT Press, 1974.
• P. E. Gill, W. Murray, and M. H. Wright, Practical Optimization, Academic Press, 1993.
• G. H. Golub and C. F. Van Loan, Matrix Computations, 2nd Edition, Johns Hopkins University Press, 1989, or
3rd edition, 1997.
• W. H. Press, B. P. Flannery, S. A. Teukolsky, and W. T. Vetterling, Numerical Recipes in C, 2nd Edition,
Cambridge University Press, 1992.
• G. Strang, Introduction to Applied Mathematics, Wellesley- Cambridge Press, 1986.
• A. E. Taylor and W. R. Mann, Advanced Calculus, 3rd Edition, John Wiley and Sons, 1983.
• L. N. Trefethen and D. Bau, III, Numerical Linear Algebra, SIAM, 1997.
• R. Weinstock, Calculus of Variations, Dover, 1974.
Chapter 2
Algebraic Linear Systems
An algebraic linear system is a set of m equations in n unknown scalars, which appear linearly. Without loss of
generality, an algebraic linear system can be written as follows:
Ax = b (2.1)
where A is an m× n matrix, x is an n-dimensional vector that collects all of the unknowns, and b is a known vector
of dimension m. In this chapter, we only consider the cases in which the entries of A, b, and x are real numbers.
Two reasons are usually offered for the importance of linear systems. The first is apparently deep, and refers
to the principle of superposition of effects. For instance, in dynamics, superposition of forces states that if force
f
1
produces acceleration a
1
(both possibly vectors) and force f
2
produces acceleration a
2
, then the combined force
f
1
+ αf
2
produces acceleration a
1
+ αa
2
. This is Newton’s second law of dynamics, although in a formulation less
common than the equivalent f = ma. Because Newton’s laws are at the basis of the entire edifice of Mechanics,
linearity appears to be a fundamental principle of Nature. However, like all physical laws, Newton’s second law is an
abstraction, and ignores viscosity, friction, turbulence, and other nonlinear effects. Linearity, then, is perhaps more in
the physicist’s mind than in reality: if nonlinear effects can be ignored, physical phenomena are linear!
A more pragmatic explanation is that linear systems are the only ones we know how to solve in general. This
argument, which is apparently more shallow than the previous one, is actually rather important. Here is why. Given
two algebraic equations in two variables,
f(x, y) = 0
g(x, y) = 0 ,
we can eliminate, say, y and obtain the equivalent system
F(x) = 0
y = h(x) .
Thus, the original system is as hard to solve as it is to find the roots of the polynomial F in a single variable. Unfortu-
nately, if f and g have degrees d
f
and d
g
, the polynomial F has generically degree d
f
d
g
.
Thus, the degree of a system of equations is, roughly speaking, the product of the degrees. For instance, a system of
m quadratic equations corresponds to a polynomial of degree 2
m
. The only case in which the exponential is harmless
is when its base is 1, that is, when the system is linear.
In this chapter, we first review a few basic facts about vectors in sections 2.1 through 2.4. More specifically, we
develop enough language to talk about linear systems and their solutions in geometric terms. In contrast with the
promise made in the introduction, these sections contain quite a few proofs. This is because a large part of the course
material is based on these notions, so we want to make sure that the foundations are sound. In addition, some of the
proofs lead to useful algorithms, and some others prove rather surprising facts. Then, in section 2.5, we characterize
the solutions of linear algebraic systems.
7
8 CHAPTER 2. ALGEBRAIC LINEAR SYSTEMS
2.1 Linear (In)dependence
Given n vectors a
1
, . . . , a
n
and n real numbers x
1
, . . . , x
n
, the vector
b =
n

j=1
x
j
a
j
(2.2)
is said to be a linear combination of a
1
, . . . , a
n
with coefficients x
1
, . . . , x
n
.
The vectors a
1
, . . . , a
n
are linearly dependent if they admit the null vector as a nonzero linear combination. In
other words, they are linearly dependent if there is a set of coefficients x
1
, . . . , x
n
, not all of which are zero, such that
n

j=1
x
j
a
j
= 0 . (2.3)
For later reference, it is useful to rewrite the last two equalities in a different form. Equation (2.2) is the same as
Ax = b (2.4)
and equation (2.3) is the same as
Ax = 0 (2.5)
where
A =
_
a
1
· · · a
n
¸
, x =
_
¸
_
x
1
.
.
.
x
n
_
¸
_ , b =
_
¸
_
b
1
.
.
.
b
m
_
¸
_ .
If you are not convinced of these equivalences, take the time to write out the components of each expression for a
small example. This is important. Make sure that you are comfortable with this.
Thus, the columns of a matrix A are dependent if there is a nonzero solution to the homogeneous system (2.5).
Vectors that are not dependent are independent.
Theorem 2.1.1 The vectors a
1
, . . . , a
n
are linearly dependent iff
1
at least one of them is a linear combination of the
others.
Proof. In one direction, dependency means that there is a nonzero vector x such that
n

j=1
x
j
a
j
= 0 .
Let x
k
be nonzero for some k. We have
n

j=1
x
j
a
j
= x
k
a
k
+
n

j=1, j=k
x
j
a
j
= 0
so that
a
k
= −
n

j=1, j=k
x
j
x
k
a
j
(2.6)
as desired. The converse is proven similarly: if
a
k
=
n

j=1, j=k
x
j
a
j
1
“iff” means “if and only if.”
2.2. BASIS 9
for some k, then
n

j=1
x
j
a
j
= 0
by letting x
k
= −1 (so that x is nonzero). ∆
2
We can make the first part of the proof above even more specific, and state the following
Lemma 2.1.2 If n nonzero vectors a
1
, . . . , a
n
are linearly dependent then at least one of them is a linear combination
of the ones that precede it.
Proof. Just let k be the last of the nonzero x
j
. Then x
j
= 0 for j > k in (2.6), which then becomes
a
k
=
n

j<k
x
j
x
k
a
j
as desired. ∆
2.2 Basis
A set a
1
, . . . , a
n
is said to be a basis for a set B of vectors if the a
j
are linearly independent and every vector in B can
be written as a linear combination of them. B is said to be a vector space if it contains all the linear combinations of
its basis vectors. In particular, this implies that every linear space contains the zero vector. The basis vectors are said
to span the vector space.
Theorem 2.2.1 Given a vector b in the vector space B and a basis a
1
, . . . , a
n
for B, the coefficients x
1
, . . . , x
n
such
that
b =
n

j=1
x
j
a
j
are uniquely determined.
Proof. Let also
b =
n

j=1
x

j
a
j
.
Then,
0 = b −b =
n

j=1
x
j
a
j

n

j=1
x

j
a
j
=
n

j=1
(x
j
−x

j
)a
j
but because the a
j
are linearly independent, this is possible only when x
j
−x

j
= 0 for every j. ∆
The previous theorem is a very important result. An equivalent formulation is the following:
If the columns a
1
, . . . , a
n
of A are linearly independent and the system Ax = b admits a solution, then
the solution is unique.
2
This symbol marks the end of a proof.
10 CHAPTER 2. ALGEBRAIC LINEAR SYSTEMS
Pause for a minute to verify that this formulation is equivalent.
Theorem 2.2.2 Two different bases for the same vector space B have the same number of vectors.
Proof. Let a
1
, . . . , a
n
and a

1
, . . . , a

n
be two different bases for B. Then each a

j
is in B (why?), and can therefore
be written as a linear combination of a
1
, . . . , a
n
. Consequently, the vectors of the set
G = a

1
, a
1
, . . . , a
n
must be linearly dependent. We call a set of vectors that contains a basis for B a generating set for B. Thus, G is a
generating set for B.
The rest of the proof now proceeds as follows: we keep removing a vectors from G and replacing them with a

vectors in such a way as to keep Ga generating set for B. Then we show that we cannot run out of a vectors before we
run out of a

vectors, which proves that n ≥ n

. We then switch the roles of a and a

vectors to conclude that n

≥ n.
This proves that n = n

.
From lemma 2.1.2, one of the vectors in G is a linear combination of those preceding it. This vector cannot be a

1
,
since it has no other vectors preceding it. So it must be one of the a
j
vectors. Removing the latter keeps Ga generating
set, since the removed vector depends on the others. Now we can add a

2
to G, writing it right after a

1
:
G = a

1
, a

2
, . . . .
G is still a generating set for B.
Let us continue this procedure until we run out of either a vectors to remove or a

vectors to add. The a vectors
cannot run out first. Suppose in fact per absurdum that G is now made only of a

vectors, and that there are still
left-over a

vectors that have not been put into G. Since the a

s form a basis, they are mutually linearly independent.
Since B is a vector space, all the a

s are in B. But then G cannot be a generating set, since the vectors in it cannot
generate the left-over a

s, which are independent of those in G. This is absurd, because at every step we have made
sure that G remains a generating set. Consequently, we must run out of a

s first (or simultaneously with the last a).
That is, n ≥ n

.
Now we can repeat the whole procedure with the roles of a vectors and a

vectors exchanged. This shows that
n

≥ n, and the two results together imply that n = n

. ∆
A consequence of this theorem is that any basis for R
m
has m vectors. In fact, the basis of elementary vectors
e
j
= jth column of the m×m identity matrix
is clearly a basis for R
m
, since any vector
b =
_
¸
_
b
1
.
.
.
b
m
_
¸
_
can be written as
b =
m

j=1
b
j
e
j
and the e
j
are clearly independent. Since this elementary basis has m vectors, theorem 2.2.2 implies that any other
basis for R
m
has m vectors.
Another consequence of theorem 2.2.2 is that n vectors of dimension m < n are bound to be dependent, since any
basis for R
m
can only have m vectors.
Since all bases for a space have the same number of vectors, it makes sense to define the dimension of a space as
the number of vectors in any of its bases.
2.3. INNER PRODUCT AND ORTHOGONALITY 11
2.3 Inner Product and Orthogonality
In this section we establish the geometric meaning of the algebraic notions of norm, inner product, projection, and
orthogonality. The fundamental geometric fact that is assumed to be known is the law of cosines: given a triangle with
sides a, b, c (see figure 2.1), we have
a
2
= b
2
+c
2
−2bc cos θ
where θ is the angle between the sides of length b and c. A special case of this law is Pythagoras’ theorem, obtained
when θ = ±π/2.
θ
c
b
a
Figure 2.1: The law of cosines states that a
2
= b
2
+c
2
−2bc cos θ.
In the previous section we saw that any vector in R
m
can be written as the linear combination
b =
m

j=1
b
j
e
j
(2.7)
of the elementary vectors that point along the coordinate axes. The length of these elementary vectors is clearly one,
because each of them goes from the origin to the unit point of one of the axes. Also, any two of these vectors form a
90-degree angle, because the coordinate axes are orthogonal by construction. How long is b? From equation (2.7) we
obtain
b = b
1
e
1
+
m

j=2
b
j
e
j
and the two vectors b
1
e
1
and

m
j=2
b
j
e
j
are orthogonal. By Pythagoras’ theorem, the square of the length b of b is
b
2
= b
2
1
+
m

j=2
b
j
e
j

2
.
Pythagoras’ theorem can now be applied again to the last sum by singling out its first term b
2
e
2
, and so forth. In
conclusion,
b
2
=
m

j=1
b
2
j
.
This result extends Pythagoras’ theorem to m dimensions.
If we define the inner product of two m-dimensional vectors as follows:
b
T
c =
m

j=1
b
j
c
j
,
then
b
2
= b
T
b . (2.8)
Thus, the squared length of a vector is the inner product of the vector with itself. Here and elsewhere, vectors are
column vectors by default, and the symbol
T
makes them into row vectors.
12 CHAPTER 2. ALGEBRAIC LINEAR SYSTEMS
Theorem 2.3.1
b
T
c = b c cos θ
where θ is the angle between b and c.
Proof. The law of cosines applied to the triangle with sides b, c, and b −c yields
b −c
2
= b
2
+c
2
−2b c cos θ
and from equation (2.8) we obtain
b
T
b + c
T
c −2b
T
c = b
T
b + c
T
c −2b c cos θ .
Canceling equal terms and dividing by -2 yields the desired result. ∆
Corollary 2.3.2 Two nonzero vectors b and c in R
m
are mutually orthogonal iff b
T
c = 0.
Proof. When θ = ±π/2, the previous theorem yields b
T
c = 0. ∆
Given two vectors b and c applied to the origin, the projection of b onto c is the vector from the origin to the point
p on the line through c that is nearest to the endpoint of b. See figure 2.2.
p
b
c
Figure 2.2: The vector from the origin to point p is the projection of b onto c. The line from the endpoint of b to p is
orthogonal to c.
Theorem 2.3.3 The projection of b onto c is the vector
p = P
c
b
where P
c
is the following square matrix:
P
c
=
cc
T
c
T
c
.
Proof. Since by definition point p is on the line through c, the projection vector p has the form p = ac, where
a is some real number. From elementary geometry, the line between p and the endpoint of b is shortest when it is
orthogonal to c:
c
T
(b −ac) = 0
2.4. ORTHOGONAL SUBSPACES AND THE RANK OF A MATRIX 13
which yields
a =
c
T
b
c
T
c
so that
p = ac = c a =
cc
T
c
T
c
b
as advertised. ∆
2.4 Orthogonal Subspaces and the Rank of a Matrix
Linear transformations map spaces into spaces. It is important to understand exactly what is being mapped into what
in order to determine whether a linear system has solutions, and if so how many. This section introduces the notion of
orthogonality between spaces, defines the null space and range of a matrix, and its rank. With these tools, we will be
able to characterize the solutions to a linear system in section 2.5. In the process, we also introduce a useful procedure
(Gram-Schmidt) for orthonormalizing a set of linearly independent vectors.
Two vector spaces A and B are said to be orthogonal to one another when every vector in A is orthogonal to every
vector in B. If vector space A is a subspace of R
m
for some m, then the orthogonal complement of A is the set of all
vectors in R
m
that are orthogonal to all the vectors in A.
Notice that complement and orthogonal complement are very different notions. For instance, the complement of
the xy plane in R
3
is all of R
3
except the xy plane, while the orthogonal complement of the xy plane is the z axis.
Theorem 2.4.1 Any basis a
1
, . . . , a
n
for a subspace A of R
m
can be extended into a basis for R
m
by adding m− n
vectors a
n+1
, . . . , a
m
.
Proof. If n = m we are done. If n < m, the given basis cannot generate all of R
m
, so there must be a vector, call
it a
n+1
, that is linearly independent of a
1
, . . . , a
n
. This argument can be repeated until the basis spans all of R
m
, that
is, until m = n. ∆
Theorem 2.4.2 (Gram-Schmidt) Given n vectors a
1
, . . . , a
n
, the following construction
r = 0
for j = 1 to n
a

j
= a
j

r
l=1
(q
T
l
a
j
)q
l
if a

j
= 0
r = r + 1
q
r
=
a

j
a

j

end
end
yields a set of orthonormal
3
vectors q
1
. . . , q
r
that span the same space as a
1
, . . . , a
n
.
Proof. We first prove by induction on r that the vectors q
r
are mutually orthonormal. If r = 1, there is little to
prove. The normalization in the above procedure ensures that q
1
has unit norm. Let us now assume that the procedure
3
Orthonormal means orthogonal and with unit norm.
14 CHAPTER 2. ALGEBRAIC LINEAR SYSTEMS
above has been performed a number j −1 of times sufficient to find r −1 vectors q
1
, . . . , q
r−1
, and that these vectors
are orthonormal (the inductive assumption). Then for any i < r we have
q
T
i
a

j
= q
T
i
a
j

r−1

l=1
(q
T
l
a
j
)q
T
i
q
l
= 0
because the term q
T
i
a
j
cancels the i-th term (q
T
i
a
j
)q
T
i
q
i
of the sum (remember that q
T
i
q
i
= 1), and the inner products
q
T
i
q
l
are zero by the inductive assumption. Because of the explicit normalization step q
r
= a

j
/a

j
, the vector q
r
, if
computed, has unit norm, and because q
T
i
a

j
= 0, it follwos that q
r
is orthogonal to all its predecessors, q
T
i
q
r
= 0 for
i = 1, . . . , r −1.
Finally, we notice that the vectors q
j
span the same space as the a
j
s, because the former are linear combinations
of the latter, are orthonormal (and therefore independent), and equal in number to the number of linearly independent
vectors in a
1
, . . . , a
n
. ∆
Theorem 2.4.3 If A is a subspace of R
m
and A

is the orthogonal complement of A in R
m
, then
dim(A) + dim(A

) = m .
Proof. Let a
1
, . . . , a
n
be a basis for A. Extend this basis to a basis a
1
, . . . , a
m
for R
m
(theorem 2.4.1). Orthonor-
malize this basis by the Gram-Schmidt procedure (theorem 2.4.2) to obtain q
1
, . . . , q
m
. By construction, q
1
, . . . , q
n
span A. Because the new basis is orthonormal, all vectors generated by q
n+1
, . . . , q
m
are orthogonal to all vectors
generated by q
1
, . . . , q
n
, so there is a space of dimension at least m− n that is orthogonal to A. On the other hand,
the dimension of this orthogonal space cannot exceed m−n, because otherwise we would have more than m vectors
in a basis for R
m
. Thus, the dimension of the orthogonal space A

is exactly m−n, as promised. ∆
We can now start to talk about matrices in terms of the subspaces associated with them. The null space null(A) of
an m× n matrix A is the space of all n-dimensional vectors that are orthogonal to the rows of A. The range of A is
the space of all m-dimensional vectors that are generated by the columns of A. Thus, x ∈ null(A) iff Ax = 0, and
b ∈ range(A) iff Ax = b for some x.
From theorem 2.4.3, if null(A) has dimension h, then the space generated by the rows of A has dimension r =
n −h, that is, A has n −h linearly independent rows. It is not obvious that the space generated by the columns of A
has also dimension r = n −h. This is the point of the following theorem.
Theorem 2.4.4 The number r of linearly independent columns of any m× n matrix A is equal to the number of its
independent rows, and
r = n −h
where h = dim(null(A)).
Proof. We have already proven that the number of independent rows is n − h. Now we show that the number of
independent columns is also n −h, by constructing a basis for range(A).
Let v
1
, . . . , v
h
be a basis for null(A), and extend this basis (theorem 2.4.1) into a basis v
1
, . . . , v
n
for R
n
. Then
we can show that the n −h vectors Av
h+1
, . . . , Av
n
are a basis for the range of A.
First, these n − h vectors generate the range of A. In fact, given an arbitrary vector b ∈ range(A), there must be
a linear combination of the columns of A that is equal to b. In symbols, there is an n-tuple x such that Ax = b. The
n-tuple x itself, being an element of R
n
, must be some linear combination of v
1
, . . . , v
n
, our basis for R
n
:
x =
n

j=1
c
j
v
j
.
2.5. THE SOLUTIONS OF A LINEAR SYSTEM 15
Thus,
b = Ax = A
n

j=1
c
j
v
j
=
n

j=1
c
j
Av
j
=
n

j=h+1
c
j
Av
j
since v
1
, . . . , v
h
span null(A), so that Av
j
= 0 for j = 1, . . . , h. This proves that the n −h vectors Av
h+1
, . . . , Av
n
generate range(A).
Second, we prove that the n − h vectors Av
h+1
, . . . , Av
n
are linearly independent. Suppose, per absurdum, that
they are not. Then there exist numbers x
h+1
, . . . , x
n
, not all zero, such that
n

j=h+1
x
j
Av
j
= 0
so that
A
n

j=h+1
x
j
v
j
= 0 .
But then the vector

n
j=h+1
x
j
v
j
is in the null space of A. Since the vectors v
1
, . . . , v
h
are a basis for null(A), there
must exist coefficients x
1
, . . . , x
h
such that
n

j=h+1
x
j
v
j
=
h

j=1
x
j
v
j
,
in conflict with the assumption that the vectors v
1
, . . . , v
n
are linearly independent. ∆
Thanks to this theorem, we can define the rank of Ato be equivalently the number of linearly independent columns
or of linearly independent rows of A:
rank(A) = dim(range(A)) = n −dim(null(A)) .
2.5 The Solutions of a Linear System
Thanks to the results of the previous sections, we now have a complete picture of the four spaces associated with an
m×n matrix A of rank r and null-space dimension h:
range(A); dimension r = rank(A)
null(A); dimension h
range(A)

; dimension m−r
null(A)

; dimension r = n −h .
The space range(A)

is called the left nullspace of the matrix, and null(A)

is called the rowspace of A. A
frequently used synonym for “range” is column space. It should be obvious from the meaning of these spaces that
null(A)

= range(A
T
)
range(A)

= null(A
T
)
where A
T
is the transpose of A, defined as the matrix obtained by exchanging the rows of A with its columns.
Theorem 2.5.1 The matrix A transforms a vector x in its null space into the zero vector, and an arbitrary vector x
into a vector in range(A).
16 CHAPTER 2. ALGEBRAIC LINEAR SYSTEMS
This allows characterizing the set of solutions to linear system as follows. Let
Ax = b
be an m×n system (m can be less than, equal to, or greater than n). Also, let
r = rank(A)
be the number of linearly independent rows or columns of A. Then,
b ∈ range(A) ⇒ no solutions
b ∈ range(A) ⇒ ∞
n−r
solutions
with the convention that ∞
0
= 1. Here, ∞
k
is the cardinality of a k-dimensional vector space.
In the first case above, there can be no linear combination of the columns (no x vector) that gives b, and the system
is said to be incompatible. In the second, compatible case, three possibilities occur, depending on the relative sizes of
r, m, n:
• When r = n = m, the system is invertible. This means that there is exactly one x that satisfies the system, since
the columns of A span all of R
n
. Notice that invertibility depends only on A, not on b.
• When r = n and m > n, the system is redundant. There are more equations than unknowns, but since b is in
the range of A there is a linear combination of the columns (a vector x) that produces b. In other words, the
equations are compatible, and exactly one solution exists.
4
• When r < n the system is underdetermined. This means that the null space is nontrivial (i.e., it has dimension
h > 0), and there is a space of dimension h = n −r of vectors x such that Ax = 0. Since b is assumed to be in
the range of A, there are solutions x to Ax = b, but then for any y ∈ null(A) also x + y is a solution:
Ax = b , Ay = 0 ⇒ A(x + y) = b
and this generates the ∞
h
= ∞
n−r
solutions mentioned above.
Notice that if r = n then n cannot possibly exceed m, so the first two cases exhaust the possibilities for r = n. Also,
r cannot exceed either m or n. All the cases are summarized in figure 2.3.
Of course, listing all possibilities does not provide an operational method for determining the type of linear system
for a given pair A, b. Gaussian elimination, and particularly its version called reduction to echelon form is such a
method, and is summarized in the next section.
2.6 Gaussian Elimination
Gaussian elimination is an important technique for solving linear systems. In addition to always yielding a solution,
no matter whether the system is invertible or not, it also allows determining the rank of a matrix.
Other solution techniques exist for linear systems. Most notably, iterative methods solve systems in a time that
depends on the accuracy required, while direct methods, like Gaussian elimination, are done in a finite amount of
time that can be bounded given only the size of a matrix. Which method to use depends on the size and structure
(e.g., sparsity) of the matrix, whether more information is required about the matrix of the system, and on numerical
considerations. More on this in chapter 3.
Consider the m×n system
Ax = b (2.9)
4
Notice that the technical meaning of “redundant” has a stronger meaning than “with more equations than unknowns.” The case r < n < mis
possible, has more equations (m) than unknowns (n), admits a solution if b ∈ range(A), but is called “underdetermined” because there are fewer
(r) independent equations than there are unknowns (see next item). Thus, “redundant” means “with exactly one solution and with more equations
than unknowns.”
2.6. GAUSSIAN ELIMINATION 17
yes no
yes no
yes no
underdetermined
redundant invertible
b in range(A)
r = n
m = n
incompatible
Figure 2.3: Types of linear systems.
which can be square or rectangular, invertible, incompatible, redundant, or underdetermined. In short, there are no
restrictions on the system. Gaussian elimination replaces the rows of this system by linear combinations of the rows
themselves until A is changed into a matrix U that is in the so-called echelon form. This means that
• Nonzero rows precede rows with all zeros. The first nonzero entry, if any, of a row, is called a pivot.
• Below each pivot is a column of zeros.
• Each pivot lies to the right of the pivot in the row above.
The same operations are applied to the rows of A and to those of b, which is transformed to a new vector c, so equality
is preserved and solving the final system yields the same solution as solving the original one.
Once the system is transformed into echelon form, we compute the solution x by backsubstitution, that is, by
solving the transformed system
Ux = c .
2.6.1 Reduction to Echelon Form
The matrix A is reduced to echelon form by a process in m − 1 steps. The first step is applied to U
(1)
= A and
c
(1)
= b. The k-th step is applied to rows k, . . . , m of U
(k)
and c
(k)
and produces U
(k+1)
and c
(k+1)
. The last step
produces U
(m)
= U and c
(m)
= c. Initially, the “pivot column index” p is set to one. Here is step k, where u
ij
denotes
entry i, j of U
(k)
:
Skip no-pivot columns If u
ip
is zero for every i = k, . . . , m, then increment p by 1. If p exceeds n stop.
5
Row exchange Now p ≤ n and u
ip
is nonzero for some k ≤ i ≤ m. Let l be one such value of i
6
. If l = k, exchange
rows l and k of U
(k)
and of c
(k)
.
Triangularization The new entry u
kp
is nonzero, and is called the pivot. For i = k + 1, . . . , m, subtract row k of
U
(k)
multiplied by u
ip
/u
kp
from row i of U
(k)
, and subtract entry k of c
(k)
multiplied by u
ip
/u
kp
from entry i
of c
(k)
. This zeros all the entries in the column below the pivot, and preserves the equality of left- and right-hand
side.
When this process is finished, U is in echelon form. In particular, if the matrix is square and if all columns have a
pivot, then U is upper-triangular.
5
“Stop” means that the entire algorithm is finished.
6
Different ways of selecting l here lead to different numerical properties of the algorithm. Selecting the largest entry in the column leads to
better round-off properties.
18 CHAPTER 2. ALGEBRAIC LINEAR SYSTEMS
2.6.2 Backsubstitution
A system
Ux = c (2.10)
in echelon form is easily solved for x. To see this, we first solve the system symbolically, leaving undetermined vari-
ables specified by their name, and then transformthis solution procedure into one that can be more readily implemented
numerically.
Let r be the index of the last nonzero row of U. Since this is the number of independent rows of U, r is the rank
of U. It is also the rank of A, because A and U admit exactly the same solutions and are equal in size. If r < m, the
last m−r equations yield a subsystem of the following form:
_
¸
_
0
.
.
.
0
_
¸
_ =
_
¸
_
c
r+1
.
.
.
c
m
_
¸
_ .
Let us call this the residual subsystem. If on the other hand r = m (obviously r cannot exceed m), there is no residual
subsystem.
If there is a residual system (i.e., r < m) and some of c
r+1
, . . . , c
m
are nonzero, then the equations corresponding
to these nonzero entries are incompatible, because they are of the form 0 = c
i
with c
i
= 0. Since no vector x can
satisfy these equations, the linear system admits no solutions: it is incompatible.
Let us now assume that either there is no residual system, or if there is one it is compatible, that is, c
r+1
= . . . =
c
m
= 0. Then, solutions exist, and they can be determined by backsubstitution, that is, by solving the equations
starting from the last one and replacing the result in the equations higher up.
Backsubstitutions works as follows. First, remove the residual system, if any. We are left with an r ×n system. In
this system, call the variables corresponding to the r columns with pivots the basic variables, and call the other n −r
the free variables. Say that the pivot columns are j
1
, . . . , j
r
. Then symbolic backsubstitution consists of the following
sequence:
for i = r downto 1
x
j
i
=
1
u
ij
i
_
_
c
i

n

l=ji+1
u
il
x
l
_
_
end
This is called symbolic backsubstitution because no numerical values are assigned to free variables. Whenever they
appear in the expressions for the basic variables, free variables are specified by name rather than by value. The final
result is a solution with as many free parameters as there are free variables. Since any value given to the free variables
leaves the equality of system (2.10) satisfied, the presence of free variables leads to an infinity of solutions.
When solving a system in echelon form numerically, however, it is inconvenient to carry around nonnumeric
symbol names (the free variables). Here is an equivalent solution procedure that makes this unnecessary. The solution
obtained by backsubstitution is an affine function
7
of the free variables, and can therefore be written in the form
x = v
0
+x
j
1
v
1
+. . . +x
j
n−r
v
n−r
(2.11)
where the x
j
i
are the free variables. The vector v
0
is the solution when all free variables are zero, and can therefore be
obtained by replacing each free variable by zero during backsubstitution. Similarly, the vector v
i
for i = 1, . . . , n −r
can be obtained by solving the homogeneous system
Ux = 0
with x
ji
= 1 and all other free variables equal to zero. In conclusion, the general solution can be obtained by running
backsubstitution n−r +1 times, once for the nonhomogeneous system, and n−r times for the homogeneous system,
with suitable values of the free variables. This yields the solution in the form (2.11).
Notice that the vectors v
1
, . . . , v
n−r
form a basis for the null space of U, and therefore of A.
7
An affine function is a linear function plus a constant.
2.6. GAUSSIAN ELIMINATION 19
2.6.3 An Example
An example will clarify both the reduction to echelon form and backsubstitution. Consider the system
Ax = b
where
U
(1)
= A =
_
_
1 3 3 2
2 6 9 5
−1 −3 3 0
_
_
, c
(1)
= b =
_
_
1
5
5
_
_
.
Reduction to echelon form transforms A and b as follows. In the first step (k = 1), there are no no-pivot columns,
so the pivot column index p stays at 1. Throughout this example, we choose a trivial pivot selection rule: we pick the
first nonzero entry at or below row k in the pivot column. For k = 1, this means that u
(1)
11
= a
11
= 1 is the pivot. In
other words, no row exchange is necessary.
8
The triangularization step subtracts row 1 multiplied by 2/1 from row 2,
and subtracts row 1 multiplied by -1/1 from row 3. When applied to both U
(1)
and c
(1)
this yields
U
(2)
=
_
_
1 3 3 2
0 0 3 1
0 0 6 2
_
_
, c
(2)
=
_
_
1
3
6
_
_
.
Notice that now (k = 2) the entries u
(2)
ip
are zero for i = 2, 3, for both p = 1 and p = 2, so p is set to 3: the second
pivot column is column 3, and u
(2)
23
is nonzero, so no row exchange is necessary. In the triangularization step, row 2
multiplied by 6/3 is subtracted from row 3 for both U
(2)
and c
(2)
to yield
U = U
(3)
=
_
_
1 3 3 2
0 0 3 1
0 0 0 0
_
_
, c = c
(3)
=
_
_
1
3
0
_
_
.
There is one zero row in the left-hand side, and the rank of U and that of A is r = 2, the number of nonzero rows.
The residual system is 0 = 0 (compatible), and r < n = 4, so the system is underdetermined, with ∞
n−r
= ∞
2
solutions.
In symbolic backsubstitution, the residual subsystem is first deleted. This yields the reduced system
_
1 3 3 2
0 0 3 1
_
x =
_
1
3
_
(2.12)
The basic variables are x
1
and x
3
, corresponding to the columns with pivots. The other two variables, x
2
and
x
4
, are free. Backsubstitution applied first to row 2 and then to row 1 yields the following expressions for the pivot
variables:
x
3
=
1
u
23
(c
2
−u
24
x
4
) =
1
3
(3 −x
4
) = 1 −
1
3
x
4
x
1
=
1
u
11
(c
1
−u
12
x
2
−u
13
x
3
−u
14
x
4
) =
1
1
(1 −3x
2
−3x
3
−2x
4
)
= 1 −3x
2
−(3 −x
4
) −2x
4
= −2 −3x
2
−x
4
so the general solution is
x =
_
¸
¸
_
−2 −3x
2
−x
4
x
2
1 −
1
3
x
4
x
4
_
¸
¸
_
=
_
¸
¸
_
−2
0
1
0
_
¸
¸
_
+x
2
_
¸
¸
_
−3
1
0
0
_
¸
¸
_
+x
4
_
¸
¸
_
−1
0

1
3
1
_
¸
¸
_
.
8
Selecting the largest entry in the column at or below row k is a frequent choice, and this would have caused rows 1 and 2 to be switched.
20 CHAPTER 2. ALGEBRAIC LINEAR SYSTEMS
This same solution can be found by the numerical backsubstitution method as follows. Solving the reduced system
(2.12) with x
2
= x
4
= 0 by numerical backsubstitution yields
x
3
=
1
3
(3 −1 · 0) = 1
x
1
=
1
1
(1 −3 · 0 −3 · 1 −2 · 0) = −2
so that
v
0
=
_
¸
¸
_
−2
0
1
0
_
¸
¸
_
.
Then v
1
is found by solving the nonzero part (first two rows) of Ux = 0 with x
2
= 1 and x
4
= 0 to obtain
x
3
=
1
3
(−1 · 0) = 0
x
1
=
1
1
(−3 · 1 −3 · 0 −2 · 0) = −3
so that
v
1
=
_
¸
¸
_
−3
1
0
0
_
¸
¸
_
.
Finally, solving the nonzero part of Ux = 0 with x
2
= 0 and x
4
= 1 leads to
x
3
=
1
3
(−1 · 1) = −
1
3
x
1
=
1
1
(−3 · 0 −3 ·
_

1
3
_
−2 · 1) = −1
so that
v
2
=
_
¸
¸
_
−1
0

1
3
1
_
¸
¸
_
and
x = v
0
+x
2
v
1
+x
4
v
2
=
_
¸
¸
_
−2
0
1
0
_
¸
¸
_
+x
2
_
¸
¸
_
−3
1
0
0
_
¸
¸
_
+x
4
_
¸
¸
_
−1
0

1
3
1
_
¸
¸
_
just as before.
As mentioned at the beginning of this section, Gaussian elimination is a direct method, in the sense that the answer
can be found in a number of steps that depends only on the size of the matrix A. In the next chapter, we study a different
2.6. GAUSSIAN ELIMINATION 21
method, based on the so-called the Singular Value Decomposition (SVD). This is an iterative method, meaning that an
exact solution usually requires an infinite number of steps, and the number of steps necessary to find an approximate
solution depends on the desired number of correct digits.
This state of affairs would seem to favor Gaussian elimination over the SVD. However, the latter yields a much
more complete answer, since it computes bases for all the four spaces mentioned above, as well as a set of quantities,
called the singular values, which provide great insight into the behavior of the linear transformation represented by
the matrix A. Singular values also allow defining a notion of approximate rank which is very useful in a large number
of applications. It also allows finding approximate solutions when the linear system in question is incompatible. In
addition, for reasons that will become apparent in the next chapter, the computation of the SVD is numerically well
behaved, much more so than Gaussian elimination. Finally, very efficient algorithms for the SVD exist. For instance,
on a regular workstation, one can compute several thousand SVDs of 5 × 5 matrices in one second. More generally,
the number of floating point operations necessary to compute the SVD of an m×n matrix is amn
2
+bn
3
where a, b
are small numbers that depend on the details of the algorithm.
22 CHAPTER 2. ALGEBRAIC LINEAR SYSTEMS
Chapter 3
The Singular Value Decomposition
In section 2, we saw that a matrix transforms vectors in its domain into vectors in its range (column space), and vectors
in its null space into the zero vector. No nonzero vector is mapped into the left null space, that is, into the orthogonal
complement of the range. In this section, we make this statement more specific by showing how unit vectors
1
in the
rowspace are transformed by matrices. This describes the action that a matrix has on the magnitudes of vectors as
well. To this end, we first need to introduce the notion of orthogonal matrices, and interpret them geometrically as
transformations between systems of orthonormal coordinates. We do this in section 3.1. Then, in section 3.2, we use
these new concepts to introduce the all-important concept of the Singular Value Decomposition (SVD). The chapter
concludes with some basic applications and examples.
3.1 Orthogonal Matrices
Let S be an n-dimensional subspace of R
m
(so that we necessarily have n ≤ m), and let v
1
, . . . , v
n
be an orthonormal
basis for S. Consider a point P in S. If the coordinates of P in R
m
are collected in an m-dimensional vector
p =
_
¸
_
p
1
.
.
.
p
m
_
¸
_ ,
and since P is in S, it must be possible to write p as a linear combination of the v
j
s. In other words, there must exist
coefficients
q =
_
¸
_
q
1
.
.
.
q
n
_
¸
_
such that
p = q
1
v
1
+. . . +q
n
v
n
= V q
where
V =
_
v
1
· · · v
n
¸
is an m×n matrix that collects the basis for S as its columns. Then for any i = 1, . . . , n we have
v
T
i
p = v
T
i
n

j=1
q
j
v
j
=
n

j=1
q
j
v
T
i
v
j
= q
i
,
since the v
j
are orthonormal. This is important, and may need emphasis:
1
Vectors with unit norm.
23
24 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION
If
p =
n

j=1
q
j
v
j
and the vectors of the basis v
1
, . . . , v
n
are orthonormal, then the coefficients q
j
are the signed mag-
nitudes of the projections of p onto the basis vectors:
q
j
= v
T
j
p . (3.1)
In matrix form,
q = V
T
p . (3.2)
Also, we can collect the n
2
equations
v
T
i
v
j
=
_
1 if i = j
0 otherwise
into the following matrix equation:
V
T
V = I (3.3)
where I is the n ×n identity matrix. A matrix V that satisfies equation (3.3) is said to be orthogonal. Thus, a matrix
is orthogonal if its columns are orthonormal. Since the left inverse of a matrix V is defined as the matrix L such that
LV = I , (3.4)
comparison with equation (3.3) shows that the left inverse of an orthogonal matrix V exists, and is equal to the
transpose of V .
Of course, this argument requires V to be full rank, so that the solution L to equation (3.4) is unique. However, V
is certainly full rank, because it is made of orthonormal columns.
Notice that V R = I cannot possibly have a solution when m > n, because the m × m identity matrix has m
linearly independent
2
columns, while the columns of V R are linear combinations of the n columns of V , so V R can
have at most n linearly independent columns.
Of course, this result is still valid when V is m×m and has orthonormal columns, since equation (3.3) still holds.
However, for square, full-rank matrices (r = m = n), the distinction between left and right inverse vanishes. In fact,
suppose that there exist matrices L and R such that LV = I and V R = I. Then L = L(V R) = (LV )R = R, so the
left and the right inverse are the same. Thus, for square orthogonal matrices, V
T
is both the left and the right inverse:
V
T
V = V V
T
= I ,
and V
T
is then simply said to be the inverse of V :
V
T
= V
−1
.
Since the matrix V V
T
contains the inner products between the rows of V (just as V
T
V is formed by the inner
products of its columns), the argument above shows that the rows of a square orthogonal matrix are orthonormal as
well. We can summarize this discussion as follows:
Theorem 3.1.1 The left inverse of an orthogonal m×n matrix V with m ≥ n exists and is equal to the transpose of
V :
V
T
V = I .
In particular, if m = n, the matrix V
−1
= V
T
is also the right inverse of V :
V square ⇒ V
−1
V = V
T
V = V V
−1
= V V
T
= I .
2
Nay, orthonormal.
3.1. ORTHOGONAL MATRICES 25
Sometimes, when m = n, the geometric interpretation of equation (3.2) causes confusion, because two interpreta-
tions of it are possible. In the interpretation given above, the point P remains the same, and the underlying reference
frame is changed from the elementary vectors e
j
(that is, from the columns of I) to the vectors v
j
(that is, to the
columns of V ). Alternatively, equation (3.2) can be seen as a transformation, in a fixed reference system, of point P
with coordinates p into a different point Q with coordinates q. This, however, is relativity, and should not be surpris-
ing: If you spin clockwise on your feet, or if you stand still and the whole universe spins counterclockwise around
you, the result is the same.
3
Consistently with either of these geometric interpretations, we have the following result:
Theorem 3.1.2 The norm of a vector x is not changed by multiplication by an orthogonal matrix V :
V x = x .
Proof.
V x
2
= x
T
V
T
V x = x
T
x = x
2
.

We conclude this section with an obvious but useful consequence of orthogonality. In section 2.3 we defined the
projection p of a vector b onto another vector c as the point on the line through c that is closest to b. This notion of
projection can be extended from lines to vector spaces by the following definition: The projection p of a point b ∈ R
n
onto a subspace C is the point in C that is closest to b.
Also, for unit vectors c, the projection matrix is cc
T
(theorem 2.3.3), and the vector b − p is orthogonal to c. An
analogous result holds for subspace projection, as the following theorem shows.
Theorem 3.1.3 Let U be an orthogonal matrix. Then the matrix UU
T
projects any vector b onto range(U). Further-
more, the difference vector between b and its projection p onto range(U) is orthogonal to range(U):
U
T
(b −p) = 0 .
Proof. A point p in range(U) is a linear combination of the columns of U:
p = Ux
where x is the vector of coefficients (as many coefficients as there are columns in U). The squared distance between b
and p is
b −p
2
= (b −p)
T
(b −p) = b
T
b + p
T
p −2b
T
p = b
T
b + x
T
U
T
Ux −2b
T
Ux .
Because of orthogonality, U
T
U is the identity matrix, so
b −p
2
= b
T
b + x
T
x −2b
T
Ux .
The derivative of this squared distance with respect to x is the vector
2x −2U
T
b
3
At least geometrically. One solution may be more efficient than the other in other ways.
26 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION
x
b
2
v
1
u σ
v
u
3
2
x
1
x
2
2
b
b
3
1
2
u σ
1 1
b
Figure 3.1: The matrix in equation (3.5) maps a circle on the plane into an ellipse in space. The two small boxes are
corresponding points.
which is zero iff
x = U
T
b ,
that is, when
p = Ux = UU
T
b
as promised.
For this value of p the difference vector b −p is orthogonal to range(U), in the sense that
U
T
(b −p) = U
T
(b −UU
T
b) = U
T
b −U
T
b = 0 .

3.2 The Singular Value Decomposition
In these notes, we have often used geometric intuition to introduce newconcepts, and we have then translated these into
algebraic statements. This approach is successful when geometry is less cumbersome than algebra, or when geometric
intuition provides a strong guiding element. The geometric picture underlying the Singular Value Decomposition is
crisp and useful, so we will use geometric intuition again. Here is the main intuition:
An m × n matrix A of rank r maps the r-dimensional unit hypersphere in rowspace(A) into an r-
dimensional hyperellipse in range(A).
This statement is stronger than saying that A maps rowspace(A) into range(A), because it also describes what
happens to the magnitudes of the vectors: a hypersphere is stretched or compressed into a hyperellipse, which is a
quadratic hypersurface that generalizes the two-dimensional notion of ellipse to an arbitrary number of dimensions. In
three dimensions, the hyperellipse is an ellipsoid, in one dimension it is a pair of points. In all cases, the hyperellipse
in question is centered at the origin.
For instance, the rank-2 matrix
A =
1

2
_
_

3

3
−3 3
1 1
_
_
(3.5)
transforms the unit circle on the plane into an ellipse embedded in three-dimensional space. Figure 3.1 shows the map
b = Ax .
3.2. THE SINGULAR VALUE DECOMPOSITION 27
Two diametrically opposite points on the unit circle are mapped into the two endpoints of the major axis of the
ellipse, and two other diametrically opposite points on the unit circle are mapped into the two endpoints of the minor
axis of the ellipse. The lines through these two pairs of points on the unit circle are always orthogonal. This result can
be generalized to any m×n matrix.
Simple and fundamental as this geometric fact may be, its proof by geometric means is cumbersome. Instead, we
will prove it algebraically by first introducing the existence of the SVD and then using the latter to prove that matrices
map hyperspheres into hyperellipses.
Theorem 3.2.1 If A is a real m×n matrix then there exist orthogonal matrices
U =
_
u
1
· · · u
m
¸
∈ R
m×m
V =
_
v
1
· · · v
n
¸
∈ R
n×n
such that
U
T
AV = Σ = diag(σ
1
, . . . , σ
p
) ∈ R
m×n
where p = min(m, n) and σ
1
≥ . . . ≥ σ
p
≥ 0. Equivalently,
A = UΣV
T
.
Proof. Let x and y be unit vectors in R
n
and R
m
, respectively, and consider the bilinear form
z = y
T
Ax .
The set
S = {x, y | x ∈ R
n
, y ∈ R
m
, x = y = 1}
is compact, so that the scalar function z(x, y) must achieve a maximum value on S, possibly at more than one point
4
.
Let u
1
, v
1
be two unit vectors in R
m
and R
n
respectively where this maximum is achieved, and let σ
1
be the corre-
sponding value of z:
max
x=y=1
y
T
Ax = u
T
1
Av
1
= σ
1
.
It is easy to see that u
1
is parallel to the vector Av
1
. If this were not the case, their inner product u
T
1
Av
1
could
be increased by rotating u
1
towards the direction of Av
1
, thereby contradicting the fact that u
T
1
Av
1
is a maximum.
Similarly, by noticing that
u
T
1
Av
1
= v
T
1
A
T
u
1
and repeating the argument above, we see that v
1
is parallel to A
T
u
1
.
By theorems 2.4.1 and 2.4.2, u
1
and v
1
can be extended into orthonormal bases for R
m
and R
n
, respectively.
Collect these orthonormal basis vectors into orthogonal matrices U
1
and V
1
. Then
U
T
1
AV
1
= S
1
=
_
σ
1
0
T
0 A
1
_
.
In fact, the first column of AV
1
is Av
1
= σ
1
u
1
, so the first entry of U
T
1
AV
1
is u
T
1
σ
1
u
1
= σ
1
, and its other entries
are u
T
j
Av
1
= 0 because Av
1
is parallel to u
1
and therefore orthogonal, by construction, to u
2
, . . . , u
m
. A similar
argument shows that the entries after the first in the first row of S
1
are zero: the row vector u
T
1
A is parallel to v
T
1
, and
therefore orthogonal to v
2
, . . . , v
n
, so that u
T
1
Av
2
= . . . = u
T
1
Av
n
= 0.
The matrix A
1
has one fewer row and column than A. We can repeat the same construction on A
1
and write
U
T
2
A
1
V
2
= S
2
=
_
σ
2
0
T
0 A
2
_
4
Actually, at least at two points: if u
T
1
Av
1
is a maximum, so is (−u
1
)
T
A(−v
1
).
28 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION
so that
_
1 0
T
0 U
T
2
_
U
T
1
AV
1
_
1 0
T
0 V
2
_
=
_
_
σ
1
0 0
T
0 σ
2
0
T
0 0 A
2
_
_
.
This procedure can be repeated until A
k
vanishes (zero rows or zero columns) to obtain
U
T
AV = Σ
where U
T
and V are orthogonal matrices obtained by multiplying together all the orthogonal matrices used in the
procedure, and
Σ = diag(σ
1
, . . . , σ
p
) .
Since matrices U and V are orthogonal, we can premultiply the matrix product in the theorem by U and postmultiply
it by V
T
to obtain
A = UΣV
T
,
which is the desired result.
It only remains to show that the elements on the diagonal of Σ are nonnegative and arranged in nonincreasing
order. To see that σ
1
≥ . . . ≥ σ
p
(where p = min(m, n)), we can observe that the successive maximization problems
that yield σ
1
, . . . , σ
p
are performed on a sequence of sets each of which contains the next. To show this, we just need
to show that σ
2
≤ σ
1
, and induction will do the rest. We have
σ
2
= max

ˆ
x=
ˆ
y=1
ˆy
T
A
1
ˆx = max

ˆ
x=
ˆ
y=1
_
0 ˆy
¸
T
S
1
_
0
ˆx
_
= max

ˆ
x=
ˆ
y=1
_
0 ˆy
¸
T
U
T
1
AV
1
_
0
ˆx
_
= max
x = y = 1
x
T
v
1
= y
T
u
1
= 0
y
T
Ax ≤ σ
1
.
To explain the last equality above, consider the vectors
x = V
1
_
0
ˆx
_
and y = U
1
_
0
ˆy
_
.
The vector x is equal to the unit vector [0 ˆx]
T
transformed by the orthogonal matrix V
1
, and is therefore itself a unit
vector. In addition, it is a linear combination of v
2
, . . . , v
n
, and is therefore orthogonal to v
1
. A similar argument
shows that y is a unit vector orthogonal to u
1
. Because x and y thus defined belong to subsets (actually sub-spheres)
of the unit spheres in R
n
and R
m
, we conclude that σ
2
≤ σ
1
.
The σ
i
are nonnegative because all these maximizations are performed on unit hyper-spheres. The σ
i
s are maxima
of the function z(x, y) which always assumes both positive and negative values on any hyper-sphere: If z(x, y) is
negative, then z(−x, y) is positive, and if x is on a hyper-sphere, so is −x. ∆
We can now review the geometric picture in figure 3.1 in light of the singular value decomposition. In the process,
we introduce some nomenclature for the three matrices in the SVD. Consider the map in figure 3.1, represented by
equation (3.5), and imagine transforming point x (the small box at x on the unit circle) into its corresponding point
b = Ax (the small box on the ellipse). This transformation can be achieved in three steps (see figure 3.2):
1. Write x in the frame of reference of the two vectors v
1
, v
2
on the unit circle that map into the major axes of the
ellipse. There are a few ways to do this, because axis endpoints come in pairs. Just pick one way, but order
v
1
, v
2
so they map into the major and the minor axis, in this order. Let us call v
1
, v
2
the two right singular
vectors of A. The corresponding axis unit vectors u
1
, u
2
on the ellipse are called left singular vectors. If we
define
V =
_
v
1
v
2
¸
,
3.2. THE SINGULAR VALUE DECOMPOSITION 29
2
x
1
x
v
2
v
1
2 2
2
v’
1
v’
2
y
y
1
u
3
y
3
u σ
2 2
u σ
1 1
σ
2 2
u’
σ
1 1
u’
x
ξ
ξ
1
ξ
η
η
1
y
η
Figure 3.2: Decomposition of the mapping in figure 3.1.
the new coordinates ξ of x become
ξ = V
T
x
because V is orthogonal.
2. Transform ξ into its image on a “straight” version of the final ellipse. “Straight” here means that the axes of the
ellipse are aligned with the y
1
, y
2
axes. Otherwise, the “straight” ellipse has the same shape as the ellipse in
figure 3.1. If the lengths of the half-axes of the ellipse are σ
1
, σ
2
(major axis first), the transformed vector η has
coordinates
η = Σξ
where
Σ =
_
_
σ
1
0
0 σ
2
0 0
_
_
is a diagonal matrix. The real, nonnegative numbers σ
1
, σ
2
are called the singular values of A.
3. Rotate the reference frame in R
m
= R
3
so that the “straight” ellipse becomes the ellipse in figure 3.1. This
rotation brings η along, and maps it to b. The components of η are the signed magnitudes of the projections of
b along the unit vectors u
1
, u
2
, u
3
that identify the axes of the ellipse and the normal to the plane of the ellipse,
so
b = Uη
where the orthogonal matrix
U =
_
u
1
u
2
u
3
¸
collects the left singular vectors of A.
30 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION
We can concatenate these three transformations to obtain
b = UΣV
T
x
or
A = UΣV
T
since this construction works for any point x on the unit circle. This is the SVD of A.
The singular value decomposition is “almost unique”. There are two sources of ambiguity. The first is in the
orientation of the singular vectors. One can flip any right singular vector, provided that the corresponding left singular
vector is flipped as well, and still obtain a valid SVD. Singular vectors must be flipped in pairs (a left vector and its
corresponding right vector) because the singular values are required to be nonnegative. This is a trivial ambiguity. If
desired, it can be removed by imposing, for instance, that the first nonzero entry of every left singular value be positive.
The second source of ambiguity is deeper. If the matrix A maps a hypersphere into another hypersphere, the axes
of the latter are not defined. For instance, the identity matrix has an infinity of SVDs, all of the form
I = UIU
T
where U is any orthogonal matrix of suitable size. More generally, whenever two or more singular values coincide,
the subspaces identified by the corresponding left and right singular vectors are unique, but any orthonormal basis can
be chosen within, say, the right subspace and yield, together with the corresponding left singular vectors, a valid SVD.
Except for these ambiguities, the SVD is unique.
Even in the general case, the singular values of a matrix A are the lengths of the semi-axes of the hyperellipse E
defined by
E = {Ax : x = 1} .
The SVD reveals a great deal about the structure of a matrix. If we define r by
σ
1
≥ . . . ≥ σ
r
> σ
r+1
= . . . = 0 ,
that is, if σ
r
is the smallest nonzero singular value of A, then
rank(A) = r
null(A) = span{v
r+1
, . . . , v
n
}
range(A) = span{u
1
, . . . , u
r
} .
The sizes of the matrices in the SVD are as follows: U is m× m, Σ is m× n, and V is n × n. Thus, Σ has the
same shape and size as A, while U and V are square. However, if m > n, the bottom (m−n) ×n block of Σ is zero,
so that the last m− n columns of U are multiplied by zero. Similarly, if m < n, the rightmost m× (n − m) block
of Σ is zero, and this multiplies the last n −m rows of V . This suggests a “small,” equivalent version of the SVD. If
p = min(m, n), we can define U
p
= U(:, 1 : p), Σ
p
= Σ(1 : p, 1 : p), and V
p
= V (:, 1 : p), and write
A = U
p
Σ
p
V
T
p
where U
p
is m×p, Σ
p
is p ×p, and V
p
is n ×p.
Moreover, if p−r singular values are zero, we can let U
r
= U(:, 1 : r), Σ
r
= Σ(1 : r, 1 : r), and V
r
= V (:, 1 : r),
then we have
A = U
r
Σ
r
V
T
r
=
r

i=1
σ
i
u
i
v
T
i
,
which is an even smaller, minimal, SVD.
Finally, both the 2-norm and the Frobenius norm
A
F
=
¸
¸
¸
_
m

i=1
n

j=1
|a
ij
|
2
3.3. THE PSEUDOINVERSE 31
and
A
2
= sup
x=0
Ax
x
are neatly characterized in terms of the SVD:
A
2
F
= σ
2
1
+. . . +σ
2
p
A
2
= σ
1
.
In the next few sections we introduce fundamental results and applications that testify to the importance of the
SVD.
3.3 The Pseudoinverse
One of the most important applications of the SVD is the solution of linear systems in the least squares sense. A linear
system of the form
Ax = b (3.6)
arising froma real-life application may or may not admit a solution, that is, a vector x that satisfies this equation exactly.
Often more measurements are available than strictly necessary, because measurements are unreliable. This leads to
more equations than unknowns (the number m of rows in A is greater than the number n of columns), and equations
are often mutually incompatible because they come from inexact measurements (incompatible linear systems were
defined in chapter 2). Even when m ≤ n the equations can be incompatible, because of errors in the measurements
that produce the entries of A. In these cases, it makes more sense to find a vector x that minimizes the norm
Ax −b
of the residual vector
r = Ax −b .
where the double bars henceforth refer to the Euclidean norm. Thus, x cannot exactly satisfy any of the m equations
in the system, but it tries to satisfy all of them as closely as possible, as measured by the sum of the squares of the
discrepancies between left- and right-hand sides of the equations.
In other circumstances, not enough measurements are available. Then, the linear system (3.6) is underdetermined,
in the sense that it has fewer independent equations than unknowns (its rank r is less than n, see again chapter 2).
Incompatibility and underdeterminacy can occur together: the system admits no solution, and the least-squares
solution is not unique. For instance, the system
x
1
+x
2
= 1
x
1
+x
2
= 3
x
3
= 2
has three unknowns, but rank 2, and its first two equations are incompatible: x
1
+ x
2
cannot be equal to both 1 and
3. A least-squares solution turns out to be x = [1 1 2]
T
with residual r = Ax − b = [1 − 1 0], which has norm

2
(admittedly, this is a rather high residual, but this is the best we can do for this problem, in the least-squares sense).
However, any other vector of the form
x

=
_
_
1
1
2
_
_

_
_
−1
1
0
_
_
is as good as x. For instance, x

= [0 2 2], obtained for α = 1, yields exactly the same residual as x (check this).
In summary, an exact solution to the system (3.6) may not exist, or may not be unique, as we learned in chapter 2.
An approximate solution, in the least-squares sense, always exists, but may fail to be unique.
32 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION
If there are several least-squares solutions, all equally good (or bad), then one of them turns out to be shorter than
all the others, that is, its norm x is smallest. One can therefore redefine what it means to “solve” a linear system so
that there is always exactly one solution. This minimum norm solution is the subject of the following theorem, which
both proves uniqueness and provides a recipe for the computation of the solution.
Theorem 3.3.1 The minimum-norm least squares solution to a linear system Ax = b, that is, the shortest vector x
that achieves the
min
x
Ax −b ,
is unique, and is given by
ˆx = V Σ

U
T
b (3.7)
where
Σ

=
_
¸
¸
¸
¸
¸
¸
¸
¸
¸
_
1/σ
1
0 · · · 0
.
.
.
1/σ
r
.
.
.
.
.
.
0
.
.
.
0 0 · · · 0
_
¸
¸
¸
¸
¸
¸
¸
¸
¸
_
is an n ×m diagonal matrix.
The matrix
A

= V Σ

U
T
is called the pseudoinverse of A.
Proof. The minimum-norm Least Squares solution to
Ax = b
is the shortest vector x that minimizes
Ax −b
that is,
UΣV
T
x −b .
This can be written as
U(ΣV
T
x −U
T
b) (3.8)
because U is an orthogonal matrix, UU
T
= I. But orthogonal matrices do not change the norm of vectors they are
applied to (theorem 3.1.2), so that the last expression above equals
ΣV
T
x −U
T
b
or, with y = V
T
x and c = U
T
b,
Σy −c .
In order to find the solution to this minimization problem, let us spell out the last expression. We want to minimize the
norm of the following vector:
_
¸
¸
¸
¸
¸
¸
¸
¸
¸
_
σ
1
0 · · · 0
0
.
.
.
· · · 0
σ
r
.
.
. 0
.
.
.
.
.
.
0 0
_
¸
¸
¸
¸
¸
¸
¸
¸
¸
_
_
¸
¸
¸
¸
¸
¸
¸
¸
_
y
1
.
.
.
y
r
y
r+1
.
.
.
y
n
_
¸
¸
¸
¸
¸
¸
¸
¸
_

_
¸
¸
¸
¸
¸
¸
¸
¸
_
c
1
.
.
.
c
r
c
r+1
.
.
.
c
m
_
¸
¸
¸
¸
¸
¸
¸
¸
_
.
3.4. LEAST-SQUARES SOLUTION OF A HOMOGENEOUS LINEAR SYSTEMS 33
The last m−r differences are of the form
0 −
_
¸
_
c
r+1
.
.
.
c
m
_
¸
_
and do not depend on the unknown y. In other words, there is nothing we can do about those differences: if some or
all the c
i
for i = r +1, . . . , m are nonzero, we will not be able to zero these differences, and each of them contributes
a residual |c
i
| to the solution. In each of the first r differences, on the other hand, the last n − r components of y are
multiplied by zeros, so they have no effect on the solution. Thus, there is freedom in their choice. Since we look for
the minimum-norm solution, that is, for the shortest vector x, we also want the shortest y, because x and y are related
by an orthogonal transformation. We therefore set y
r+1
= . . . = y
n
= 0. In summary, the desired y has the following
components:
y
i
=
c
i
σ
i
for i = 1, . . . , r
y
i
= 0 for i = r + 1, . . . , n .
When written as a function of the vector c, this is
y = Σ
+
c .
Notice that there is no other choice for y, which is therefore unique: minimum residual forces the choice of y
1
, . . . , y
r
,
and minimum-norm solution forces the other entries of y. Thus, the minimum-norm, least-squares solution to the
original system is the unique vector
ˆx = V y = V Σ
+
c = V Σ
+
U
T
b
as promised. The residual, that is, the norm of Ax − b when x is the solution vector, is the norm of Σy − c, since
this vector is related to Ax −b by an orthogonal transformation (see equation (3.8)). In conclusion, the square of the
residual is
Ax −b
2
= Σy −c
2
=
m

i=r+1
c
2
i
=
m

i=r+1
(u
T
i
b)
2
which is the projection of the right-hand side vector b onto the complement of the range of A. ∆
3.4 Least-Squares Solution of a Homogeneous Linear Systems
Theorem 3.3.1 works regardless of the value of the right-hand side vector b. When b = 0, that is, when the system is
homogeneous, the solution is trivial: the minimum-norm solution to
Ax = 0 (3.9)
is
x = 0 ,
which happens to be an exact solution. Of course it is not necessarily the only one (any vector in the null space of A
is also a solution, by definition), but it is obviously the one with the smallest norm.
Thus, x = 0 is the minimum-norm solution to any homogeneous linear system. Although correct, this solution is
not too interesting. In many applications, what is desired is a nonzero vector x that satisfies the system (3.9) as well
as possible. Without any constraints on x, we would fall back to x = 0 again. For homogeneous linear systems, the
meaning of a least-squares solution is therefore usually modified, once more, by imposing the constraint
x = 1
34 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION
on the solution. Unfortunately, the resulting constrained minimization problem does not necessarily admit a unique
solution. The following theorem provides a recipe for finding this solution, and shows that there is in general a whole
hypersphere of solutions.
Theorem 3.4.1 Let
A = UΣV
T
be the singular value decomposition of A. Furthermore, let v
n−k+1
, . . . , v
n
be the k columns of V whose correspond-
ing singular values are equal to the last singular value σ
n
, that is, let k be the largest integer such that
σ
n−k+1
= . . . = σ
n
.
Then, all vectors of the form
x = α
1
v
n−k+1
+. . . +α
k
v
n
(3.10)
with
α
2
1
+. . . +α
2
k
= 1 (3.11)
are unit-norm least squares solutions to the homogeneous linear system
Ax = 0,
that is, they achieve the
min
x=1
Ax .
Note: when σ
n
is greater than zero the most common case is k = 1, since it is very unlikely that different singular
values have exactly the same numerical value. When A is rank deficient, on the other case, it may often have more
than one singular value equal to zero. In any event, if k = 1, then the minimum-norm solution is unique, x = v
n
. If
k > 1, the theorem above shows how to express all solutions as a linear combination of the last k columns of V .
Proof. The reasoning is very similar to that for the previous theorem. The unit-norm Least Squares solution to
Ax = 0
is the vector x with x = 1 that minimizes
Ax
that is,
UΣV
T
x .
Since orthogonal matrices do not change the norm of vectors they are applied to (theorem 3.1.2), this norm is the same
as
ΣV
T
x
or, with y = V
T
x,
Σy .
Since V is orthogonal, x = 1 translates to y = 1. We thus look for the unit-norm vector y that minimizes the
norm (squared) of Σy, that is,
σ
2
1
y
2
1
+. . . +σ
2
n
y
2
n
.
This is obviously achieved by concentrating all the (unit) mass of y where the σs are smallest, that is by letting
y
1
= . . . = y
n−k
= 0. (3.12)
From y = V
T
x we obtain x = V y = y
1
v
1
+. . . +y
n
v
n
, so that equation (3.12) is equivalent to equation (3.10) with
α
1
= y
n−k+1
, . . . , α
k
= y
n
, and the unit-norm constraint on y yields equation (3.11). ∆
Section 3.5 shows a sample use of theorem 3.4.1.
3.5. SVD LINE FITTING 35
3.5 SVD Line Fitting
The Singular Value Decomposition of a matrix yields a simple method for fitting a line to a set of points on the plane.
3.5.1 Fitting a Line to a Set of Points
Let p
i
= (x
i
, y
i
)
T
be a set of m ≥ 2 points on the plane, and let
ax +by −c = 0
be the equation of a line. If the lefthand side of this equation is multiplied by a nonzero constant, the line does not
change. Thus, we can assume without loss of generality that
n = a
2
+b
2
= 1 , (3.13)
where the unit vector n = (a, b)
T
, orthogonal to the line, is called the line normal.
The distance from the line to the origin is |c| (see figure 3.3), and the distance between the line n and a point p
i
is
equal to
d
i
= |ax
i
+by
i
−c| = |p
T
i
n −c| . (3.14)
p
i
a
b
|c|
Figure 3.3: The distance between point p
i
= (x
i
, y
i
)
T
and line ax +by −c = 0 is |ax
i
+by
i
−c|.
The best-fit line minimizes the sum of the squared distances. Thus, if we let d = (d
1
, . . . , d
m
) and P =
(p
1
. . . , p
m
)
T
, the best-fit line achieves the
min
n=1
d
2
= min
n=1
Pn −c1
2
. (3.15)
In equation (3.15), 1 is a vector of m ones.
3.5.2 The Best Line Fit
Since the third line parameter c does not appear in the constraint (3.13), at the minimum (3.15) we must have
∂d
2
∂c
= 0 . (3.16)
If we define the centroid p of all the points p
i
as
p =
1
m
P
T
1 ,
36 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION
equation (3.16) yields
∂d
2
∂c
=

∂c
_
n
T
P
T
−c1
T
_
(Pn −1c)
=

∂c
_
n
T
P
T
Pn +c
2
1
T
1 −2n
T
P
T
c1
_
= 2
_
mc −n
T
P
T
1
_
= 0
from which we obtain
c =
1
m
n
T
P
T
1 ,
that is,
c = p
T
n .
By replacing this expression into equation (3.15), we obtain
min
n=1
d
2
= min
n=1
Pn −1p
T
n
2
= min
n=1
Qn
2
,
where Q = P − 1p
T
collects the centered coordinates of the m points. We can solve this constrained minimization
problem by theorem 3.4.1. Equivalently, and in order to emphasize the geometric meaning of signular values and
vectors, we can recall that if n is on a circle, the shortest vector of the form Qn is obtained when n is the right singular
vector v
2
corresponding to the smaller σ
2
of the two singular values of Q. Furthermore, since Qv
2
has norm σ
2
, the
residue is
min
n=1
d = σ
2
and more specifically the distances d
i
are given by
d = σ
2
u
2
where u
2
is the left singular vector corresponding to σ
2
. In fact, when n = v
2
, the SVD
Q = UΣV
T
=
2

i=1
σ
i
u
i
v
T
i
yields
Qn = Qv
2
=
2

i=1
σ
i
u
i
v
T
i
v
2
= σ
2
u
2
because v
1
and v
2
are orthonormal vectors.
To summarize, to fit a line (a, b, c) to a set of m points p
i
collected in the m × 2 matrix P = (p
1
. . . , p
m
)
T
,
proceed as follows:
1. compute the centroid of the points (1 is a vector of m ones):
p =
1
m
P
T
1
2. form the matrix of centered coordinates:
Q = P −1p
T
3. compute the SVD of Q:
Q = UΣV
T
3.5. SVD LINE FITTING 37
4. the line normal is the second column of the 2 ×2 matrix V :
n = (a, b)
T
= v
2
,
5. the third coefficient of the line is
c = p
T
n
6. the residue of the fit is
min
n=1
d = σ
2
The following matlab code implements the line fitting method.
function [l, residue] = linefit(P)
% check input matrix sizes
[m n] = size(P);
if n ˜= 2, error(’matrix P must be m x 2’), end
if m < 2, error(’Need at least two points’), end
one = ones(m, 1);
% centroid of all the points
p = (P’ * one) / m;
% matrix of centered coordinates
Q = P - one * p’;
[U Sigma V] = svd(Q);
% the line normal is the second column of V
n = V(:, 2);
% assemble the three line coefficients into a column vector
l = [n ; p’ * n];
% the smallest singular value of Q
% measures the residual fitting error
residue = Sigma(2, 2);
A useful exercise is to think how this procedure, or something close to it, can be adapted to fit a set of data points
in R
m
with an affine subspace of given dimension n. An affine subspace is a linear subspace plus a point, just like an
arbitrary line is a line through the origin plus a point. Here “plus” means the following. Let L be a linear space. Then
an affine space has the form
A = p +L = {a | a = p + l and l ∈ L} .
Hint: minimizing the distance between a point and a subspace is equivalent to maximizing the norm of the projection
of the point onto the subspace. The fitting problem (including fitting a line to a set of points) can be cast either as a
maximization or a minimization problem.
38 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION
Chapter 4
Function Optimization
There are three main reasons why most problems in robotics, vision, and arguably every other science or endeavor
take on the form of optimization problems. One is that the desired goal may not be achievable, and so we try to get as
close as possible to it. The second reason is that there may be more ways to achieve the goal, and so we can choose
one by assigning a quality to all the solutions and selecting the best one. The third reason is that we may not know
how to solve the system of equations f(x) = 0, so instead we minimize the norm f(x), which is a scalar function of
the unknown vector x.
We have encountered the first two situations when talking about linear systems. The case in which a linear system
admits exactly one exact solution is simple but rare. More often, the system at hand is either incompatible (some say
overconstrained) or, at the opposite end, underdetermined. In fact, some problems are both, in a sense. While these
problems admit no exact solution, they often admit a multitude of approximate solutions. In addition, many problems
lead to nonlinear equations.
Consider, for instance, the problem of Structure From Motion (SFM) in computer vision. Nonlinear equations
describe how points in the world project onto the images taken by cameras at given positions in space. Structure from
motion goes the other way around, and attempts to solve these equations: image points are given, and one wants to
determine where the points in the world and the cameras are. Because image points come from noisy measurements,
they are not exact, and the resulting system is usually incompatible. SFM is then cast as an optimization problem.
On the other hand, the exact system (the one with perfect coefficients) is often close to being underdetermined. For
instance, the images may be insufficient to recover a certain shape under a certain motion. Then, an additional criterion
must be added to define what a “good” solution is. In these cases, the noisy system admits no exact solutions, but has
many approximate ones.
The term “optimization” is meant to subsume both minimization and maximization. However, maximizing the
scalar function f(x) is the same as minimizing −f(x), so we consider optimization and minimization to be essentially
synonyms. Usually, one is after global minima. However, global minima are hard to find, since they involve a universal
quantifier: x

is a global minimum of f if for every other x we have f(x) ≥ f(x

). Global minization techniques
like simulated annealing have been proposed, but their convergence properties depend very strongly on the problem at
hand. In this chapter, we consider local minimization: we pick a starting point x
0
, and we descend in the landscape of
f(x) until we cannot go down any further. The bottom of the valley is a local minimum.
Local minimization is appropriate if we know how to pick an x
0
that is close to x

. This occurs frequently in
feedback systems. In these systems, we start at a local (or even a global) minimum. The system then evolves and
escapes from the minimum. As soon as this occurs, a control signal is generated to bring the system back to the
minimum. Because of this immediate reaction, the old minimum can often be used as a starting point x
0
when looking
for the new minimum, that is, when computing the required control signal. More formally, we reach the correct
minimum x

as long as the initial point x
0
is in the basin of attraction of x

, defined as the largest neighborhood of x

in which f(x) is convex.
Good references for the discussion in this chapter are Matrix Computations, Practical Optimization, and Numerical
Recipes in C, all of which are listed with full citations in section 1.4.
39
40 CHAPTER 4. FUNCTION OPTIMIZATION
4.1 Local Minimization and Steepest Descent
Suppose that we want to find a local minimum for the scalar function f of the vector variable x, starting from an initial
point x
0
. Picking an appropriate x
0
is crucial, but also very problem-dependent. We start from x
0
, and we go downhill.
At every step of the way, we must make the following decisions:
• Whether to stop.
• In what direction to proceed.
• How long a step to take.
In fact, most minimization algorithms have the following structure:
k = 0
while x
k
is not a minimum
compute step direction p
k
with p
k
= 1
compute step size α
k
x
k+1
= x
k

k
p
k
k = k + 1
end.
Different algorithms differ in how each of these instructions is performed.
It is intuitively clear that the choice of the step size α
k
is important. Too small a step leads to slow convergence,
or even to lack of convergence altogether. Too large a step causes overshooting, that is, leaping past the solution. The
most disastrous consequence of this is that we may leave the basin of attraction, or that we oscillate back and forth
with increasing amplitudes, leading to instability. Even when oscillations decrease, they can slow down convergence
considerably.
What is less obvious is that the best direction of descent is not necessarily, and in fact is quite rarely, the direction
of steepest descent, as we now show. Consider a simple but important case,
f(x) = c + a
T
x +
1
2
x
T
Qx (4.1)
where Q is a symmetric, positive definite matrix. Positive definite means that for every nonzero x the quantity x
T
Qx
is positive. In this case, the graph of f(x) −c is a plane a
T
x plus a paraboloid.
Of course, if f were this simple, no descent methods would be necessary. In fact the minimum of f can be found
by setting its gradient to zero:
∂f
∂x
= a +Qx = 0
so that the minimum x

is the solution to the linear system
Qx = −a . (4.2)
Since Q is positive definite, it is also invertible (why?), and the solution x

is unique. However, understanding the
behavior of minimization algorithms in this simple case is crucial in order to establish the convergence properties of
these algorithms for more general functions. In fact, all smooth functions can be approximated by paraboloids in a
sufficiently small neighborhood of any point.
Let us therefore assume that we minimize f as given in equation (4.1), and that at every step we choose the
direction of steepest descent. In order to simplify the mathematics, we observe that if we let
˜ e(x) =
1
2
(x −x

)
T
Q(x −x

)
then we have
˜ e(x) = f(x) −c +
1
2
x
∗T
Qx

= f(x) −f(x

) (4.3)
4.1. LOCAL MINIMIZATION AND STEEPEST DESCENT 41
so that ˜ e and f differ only by a constant. In fact,
˜ e(x) =
1
2
(x
T
Qx + x
∗T
Qx

−2x
T
Qx

) =
1
2
x
T
Qx + a
T
x +
1
2
x
∗T
Qx

= f(x) −c +
1
2
x
∗T
Qx

and from equation (4.2) we obtain
f(x

) = c + a
T
x

+
1
2
x
∗T
Qx

= c −x
∗T
Qx

+
1
2
x
∗T
Qx

= c −
1
2
x
∗T
Qx

.
The result (4.3),
˜ e(x) = f(x) −f(x

) ,
is rather interesting in itself. It says that adding a linear term a
T
x (and a constant c) to a paraboloid
1
2
x
T
Qx merely
shifts the bottom of the paraboloid, both in position (x

rather than 0) and value (c −
1
2
x
∗T
Qx

rather than zero).
Adding the linear term does not “warp” or “tilt” the shape of the paraboloid in any way.
Since ˜ e is simpler, we consider that we are minimizing ˜ e rather than f. In addition, we can let
y = x −x

,
that is, we can shift the origin of the domain to x

, and study the function
e(y) =
1
2
y
T
Qy
instead of f or ˜ e, without loss of generality. We will transform everything back to f and x once we are done. Of
course, by construction, the new minimum is at
y

= 0
where e reaches a value of zero:
e(y

) = e(0) = 0 .
However, we let our steepest descent algorithm find this minimum by starting from the initial point
y
0
= x
0
−x

.
At every iteration k, the algorithm chooses the direction of steepest descent, which is in the direction
p
k
= −
g
k
g
k

opposite to the gradient of e evaluated at y
k
:
g
k
= g(y
k
) =
∂e
∂y
¸
¸
¸
¸
y=y
k
= Qy
k
.
We select for the algorithm the most favorable step size, that is, the one that takes us from y
k
to the lowest point in
the direction of p
k
. This can be found by differentiating the function
e(y
k
+αp
k
) =
1
2
(y
k
+αp
k
)
T
Q(y
k
+αp
k
)
with respect to α, and setting the derivative to zero to obtain the optimal step α
k
. We have
∂e(y
k
+αp
k
)
∂α
= (y
k
+αp
k
)
T
Qp
k
and setting this to zero yields
α
k
= −
(Qy
k
)
T
p
k
p
T
k
Qp
k
= −
g
T
k
p
k
p
T
k
Qp
k
= g
k

p
T
k
p
k
p
T
k
Qp
k
= g
k

g
T
k
g
k
g
T
k
Qg
k
. (4.4)
42 CHAPTER 4. FUNCTION OPTIMIZATION
Thus, the basic step of our steepest descent can be written as follows:
y
k+1
= y
k
+g
k

g
T
k
g
k
g
T
k
Qg
k
p
k
that is,
y
k+1
= y
k

g
T
k
g
k
g
T
k
Qg
k
g
k
. (4.5)
How much closer did this step bring us to the solution y

= 0? In other words, how much smaller is e(y
k+1
),
relative to the value e(y
k
) at the previous step? The answer is, often not much, as we shall now prove. The arguments
and proofs below are adapted from D. G. Luenberger, Introduction to Linear and Nonlinear Programming, Addison-
Wesley, 1973.
From the definition of e and from equation (4.5) we obtain
e(y
k
) −e(y
k+1
)
e(y
k
)
=
y
T
k
Qy
k
−y
T
k+1
Qy
k+1
y
T
k
Qy
k
=
y
T
k
Qy
k

_
y
k

g
T
k
g
k
g
T
k
Qg
k
g
k
_
T
Q
_
y
k

g
T
k
g
k
g
T
k
Qg
k
g
k
_
y
T
k
Qy
k
=
2
g
T
k
g
k
g
T
k
Qg
k
g
T
k
Qy
k

_
g
T
k
g
k
g
T
k
Qg
k
_
2
g
T
k
Qg
k
y
T
k
Qy
k
=
2g
T
k
g
k
g
T
k
Qy
k
−(g
T
k
g
k
)
2
y
T
k
Qy
k
g
T
k
Qg
k
.
Since Q is invertible we have
g
k
= Qy
k
⇒ y
k
= Q
−1
g
k
and
y
T
k
Qy
k
= g
T
k
Q
−1
g
k
so that
e(y
k
) −e(y
k+1
)
e(y
k
)
=
(g
T
k
g
k
)
2
g
T
k
Q
−1
g
k
g
T
k
Qg
k
.
This can be rewritten as follows by rearranging terms:
e(y
k+1
) =
_
1 −
(g
T
k
g
k
)
2
g
T
k
Q
−1
g
k
g
T
k
Qg
k
_
e(y
k
) (4.6)
so if we can bound the expression in parentheses we have a bound on the rate of convergence of steepest descent. To
this end, we introduce the following result.
Lemma 4.1.1 (Kantorovich inequality) Let Qbe a positive definite, symmetric, n×n matrix. For any vector y there
holds
(y
T
y)
2
y
T
Q
−1
y y
T
Qy


1
σ
n

1

n
)
2
where σ
1
and σ
n
are, respectively, the largest and smallest singular values of Q.
4.1. LOCAL MINIMIZATION AND STEEPEST DESCENT 43
Proof. Let
Q = UΣU
T
be the singular value decomposition of the symmetric (hence V = U) matrix Q. Because Q is positive definite, all its
singular values are strictly positive, since the smallest of them satisfies
σ
n
= min
y=1
y
T
Qy > 0
by the definition of positive definiteness. If we let
z = U
T
y
we have
(y
T
y)
2
y
T
Q
−1
y y
T
Qy
=
(y
T
U
T
Uy)
2
y
T

−1
U
T
y y
T
UΣU
T
y
=
(z
T
z)
2
z
T
Σ
−1
z z
T
Σz
=
1/

n
i=1
θ
i
σ
i

n
i=1
θ
i

i
=
φ(σ)
ψ(σ)
(4.7)
where the coefficients
θ
i
=
z
2
i
z
2
add up to one. If we let
σ =
n

i=1
θ
i
σ
i
, (4.8)
then the numerator φ(σ) in (4.7) is 1/σ. Of course, there are many ways to choose the coefficients θ
i
to obtain a
particular value of σ. However, each of the singular values σ
j
can be obtained by letting θ
j
= 1 and all other θ
i
to
zero. Thus, the values 1/σ
j
for j = 1, . . . , n are all on the curve 1/σ. The denominator ψ(σ) in (4.7) is a convex
combination of points on this curve. Since 1/σ is a convex function of σ, the values of the denominator ψ(σ) of (4.7)
must be in the shaded area in figure 4.1. This area is delimited from above by the straight line that connects point

1
, 1/σ
1
) with point (σ
n
, 1/σ
n
), that is, by the line with ordinate
λ(σ) = (σ
1

n
−σ)/(σ
1
σ
n
) .
For the same vector of coefficients θ
i
, the values of φ(σ), ψ(σ), and λ(σ) are on the vertical line corresponding to
the value of σ given by (4.8). Thus an appropriate bound is
φ(σ)
ψ(σ)
≥ min
σ
1
≤σ≤σ
n
φ(σ)
λ(σ)
= min
σ
1
≤σ≤σ
n
1/σ

1

n
−σ)/(σ
1
σ
n
)
.
The minimum is achieved at σ = (σ
1

n
)/2, yielding the desired result. ∆
Thanks to this lemma, we can state the main result on the convergence of the method of steepest descent.
Theorem 4.1.2 Let
f(x) = c + a
T
x +
1
2
x
T
Qx
be a quadratic function of x, with Q symmetric and positive definite. For any x
0
, the method of steepest descent
x
k+1
= x
k

g
T
k
g
k
g
T
k
Qg
k
g
k
(4.9)
where
g
k
= g(x
k
) =
∂f
∂x
¸
¸
¸
¸
x=x
k
= a +Qx
k
44 CHAPTER 4. FUNCTION OPTIMIZATION
σ
2 1
σ σ σ
n
λ(σ)
φ(σ)
ψ(σ)
φ,ψ,λ
σ
Figure 4.1: Kantorovich inequality.
converges to the unique minimum point
x

= −Q
−1
a
of f. Furthermore, at every step k there holds
f(x
k+1
) −f(x

) ≤
_
σ
1
−σ
n
σ
1

n
_
2
(f(x
k
) −f(x

))
where σ
1
and σ
n
are, respectively, the largest and smallest singular value of Q.
Proof. From the definitions
y = x −x

and e(y) =
1
2
y
T
Qy (4.10)
we immediately obtain the expression for steepest descent in terms of f and x. By equations (4.3) and (4.6) and the
Kantorovich inequality we obtain
f(x
k+1
) −f(x

) = e(y
k+1
) =
_
1 −
(g
T
k
g
k
)
2
g
T
k
Q
−1
g
k
g
T
k
Qg
k
_
e(y
k
) ≤
_
1 −

1
σ
n

1

n
)
2
_
e(y
k
) (4.11)
=
_
σ
1
−σ
n
σ
1

n
_
2
(f(x
k
) −f(x

)) . (4.12)
Since the ratio in the last term is smaller than one, it follows immediately that f(x
k
) − f(x

) → 0 and hence, since
the minimum of f is unique, that x
k
→x

. ∆
The ratio κ(Q) = σ
1

n
is called the condition number of Q. The larger the condition number, the closer the
fraction (σ
1
− σ
n
)/(σ
1
+ σ
n
) is to unity, and the slower convergence. It is easily seen why this happens in the case
4.1. LOCAL MINIMIZATION AND STEEPEST DESCENT 45
p
0
x
*
x
0 1
p
Figure 4.2: Trajectory of steepest descent.
in which x is a two-dimensional vector, as in figure 4.2, which shows the trajectory x
k
superimposed on a set of
isocontours of f(x).
There is one good, but very precarious case, namely, when the starting point x
0
is at one apex (tip of either axis)
of an isocontour ellipse. In that case, one iteration will lead to the minimum x

. In all other cases, the line in the
direction p
k
of steepest descent, which is orthogonal to the isocontour at x
k
, will not pass through x

. The minimum
of f along that line is tangent to some other, lower isocontour. The next step is orthogonal to the latter isocontour (that
is, parallel to the gradient). Thus, at every step the steepest descent trajectory is forced to make a ninety-degree turn.
If isocontours were circles (σ
1
= σ
n
) centered at x

, then the first turn would make the new direction point to x

, and
minimization would get there in just one more step. This case, in which κ(Q) = 1, is consistent with our analysis,
because then
σ
1
−σ
n
σ
1

n
= 0 .
The more elongated the isocontours, that is, the greater the condition number κ(Q), the farther away a line orthogonal
to an isocontour passes from x

, and the more steps are required for convergence.
For general (that is, non-quadratic) f, the analysis above applies once x
k
gets close enough to the minimum, so
that f is well approximated by a paraboloid. In this case, Q is the matrix of second derivatives of f with respect to x,
and is called the Hessian of f. In summary, steepest descent is good for functions that have a well conditioned Hessian
near the minimum, but can become arbitrarily slow for poorly conditioned Hessians.
To characterize the speed of convergence of different minimization algorithms, we introduce the notion of the
order of convergence. This is defined as the largest value of q for which the
lim
k→∞
x
k+1
−x

x
k
−x

q
is finite. If β is this limit, then close to the solution (that is, for large values of k) we have
x
k+1
−x

≈ βx
k
−x

q
for a minimization method of order q. In other words, the distance of x
k
from x

is reduced by the q-th power at every
step, so the higher the order of convergence, the better. Theorem 4.1.2 implies that steepest descent has at best a linear
order of convergence. In fact, the residuals |f(x
k
) − f(x

)| in the values of the function being minimized converge
46 CHAPTER 4. FUNCTION OPTIMIZATION
linearly. Since the gradient of f approaches zero when x
k
tends to x

, the arguments x
k
to f can converge to x

even
more slowly.
To complete the steepest descent algorithm we need to specify how to check whether a minimum has been reached.
One criterion is to check whether the value of f(x
k
) has significantly decreased from f(x
k−1
). Another is to check
whether x
k
is significantly different from x
k−1
. Close to the minimum, the derivatives of f are close to zero, so
|f(x
k
) − f(x
k−1
)| may be very small but x
k
− x
k−1
may still be relatively large. Thus, the check on x
k
is more
stringent, and therefore preferable in most cases. In fact, usually one is interested in the value of x

, rather than in that
of f(x

). In summary, the steepest descent algorithm can be stopped when
x
k
−x
k−1
<
where the positive constant is provided by the user.
In our analysis of steepest descent, we used the Hessian Qin order to compute the optimal step size α (see equation
(4.4)). We used Qbecause it was available, but its computation during steepest descent would in general be overkill. In
fact, only gradient information is necessary to find p
k
, and a line search in the direction of p
k
can be used to determine
the step size α
k
. In contrast, the Hessian of f(x) requires computing
_
n
2
_
second derivatives if x is an n-dimensional
vector.
Using line search to find α
k
guarantees that a minimum in the direction p
k
is actually reached even when the
parabolic approximation is inadequate. Here is how line search works.
Let
h(α) = f(x
k
+αp
k
) (4.13)
be the scalar function of one variable that is obtained by restricting the function f to the line through the current point
x
k
and in the direction of p
k
. Line search first determines two points a, c that bracket the desired minimum α
k
, in the
sense that a ≤ α
k
≤ c, and then picks a point between a and c, say, b = (a + c)/2. The only difficulty here is to
find c. In fact, we can set a = 0, corresponding through equation (4.13) to the starting point x
k
. A point c that is on
the opposite side of the minimum with respect to a can be found by increasing α through values α
1
= a, α
2
, . . . until
h(α
i
) is greater than h(α
i−1
). Then, if we can assume that h is convex between α
1
and α
i
, we can set c = α
i
. In
fact, the derivative of h at a is negative, so the function is initially decreasing, but it is increasing between α
i−1
and
α
i
= c, so the minimum must be somewhere between a and c. Of course, if we cannot assume convexity, we may find
the wrong minimum, but there is no general-purpose fix to this problem.
Line search now proceeds by shrinking the bracketing triple (a, b, c) until c−a is smaller than the desired accuracy
in determining α
k
. Shrinking works as follows:
if b −a > c −b
u = (a +b)/2
if f(u) > f(b)
(a, b, c) = (u, b, c)
otherwise
(a, b, c) = (a, u, b)
end
otherwise
u = (b +c)/2
if f(u) > f(b)
(a, b, c) = (a, b, u)
otherwise
(a, b, c) = (b, u, c)
end
end.
4.2. NEWTON’S METHOD 47
It is easy to see that in each case the bracketing triple (a, b, c) preserves the property that f(b) ≤ f(a) and
f(b) ≤ f(c), and therefore the minimum is somewhere between a and c. In addition, at every step the interval (a, c)
shrinks to 3/4 of its previous size, so line search will find the minimum in a number of steps that is logarithmic in the
desired accuracy.
4.2 Newton’s Method
If a function can be well approximated by a paraboloid in the region in which minimization is performed, the analysis
in the previous section suggests a straight-forward fix to the slow convergence of steepest descent. In fact, equation
(4.2) tells us how to jump in one step from the starting point x
0
to the minimum x

. Of course, when f(x) is not
exactly a paraboloid, the new value x
1
will be different from x

. Consequently, iterations are needed, but convergence
can be expected to be faster. This is the idea of Newton’s method, which we now summarize. Let
f(x
k
+ ∆x) ≈ f(x
k
) + g
T
k
∆x +
1
2
∆x
T
Q
k
∆x (4.14)
be the first terms of the Taylor series expansion of f about the current point x
k
, where
g
k
= g(x
k
) =
∂f
∂x
¸
¸
¸
¸
x=x
k
and
Q
k
= Q(x
k
) =

2
f
∂x∂x
T
¸
¸
¸
¸
x=x
k
=
_
¸
¸
_

2
f
∂x
2
1
· · ·

2
f
∂x1∂xn
.
.
.
.
.
.

2
f
∂xn∂x1
· · ·

2
f
∂x
2
n
_
¸
¸
_
x=x
k
are the gradient and Hessian of f evaluated at the current point x
k
. Notice that even when f is a paraboloid, the
gradient g
k
is different from a as used in equation (4.1). In fact, a and Q are the coefficients of the Taylor expansion
of f around point x = 0, while g
k
and Q
k
are the coefficients of the Taylor expansion of f around the current point
x
k
. In other words, gradient and Hessian are constantly reevaluated in Newton’s method.
To the extent that approximation (4.14) is valid, we can set the derivatives of f(x
k
+ ∆x) with respect to ∆x to
zero, and obtain, analogously to equation (4.2), the linear system
Q
k
∆x = −g
k
, (4.15)
whose solution ∆x
k
= α
k
p
k
yields at the same time the step direction p
k
= ∆x
k
/∆x
k
and the step size α
k
=
∆x
k
. The direction is of course undefined once the algorithm has reached a minimum, that is, when α
k
= 0.
A minimization algorithm in which the step direction p
k
and size α
k
are defined in this manner is called Newton’s
method. The corresponding p
k
is termed the Newton direction, and the step defined by equation (4.15) is the Newton
step.
The greater speed of Newton’s method over steepest descent is borne out by analysis: while steepest descent has a
linear order of convergence, Newton’s method is quadratic. In fact, let
y(x) = x −Q(x)
−1
g(x)
be the place reached by a Newton step starting at x (see equation (4.15)), and suppose that at the minimum x

the
Hessian Q(x

) is nonsingular. Then
y(x

) = x

because g(x

) = 0, and
x
k+1
−x

= y(x
k
) −x

= y(x
k
) −y(x

) .
48 CHAPTER 4. FUNCTION OPTIMIZATION
From the mean-value theorem, we have
x
k+1
−x

= y(x
k
) −y(x

) ≤
_
_
_
_
_
∂y
∂x
T
_
x=x

(x
k
−x

)
_
_
_
_
+
1
2
¸
¸
¸
¸

2
y
∂x∂x
T
¸
¸
¸
¸
x=
ˆ
x
x
k
−x

2
where ˆx is some point on the line between x

and x
k
. Since y(x

) = x

, the first derivatives of y at x

are zero, so that
the first term in the right-hand side above vanishes, and
x
k+1
−x

≤ c x
k
−x

2
where c depends on third-order derivatives of f near x

. Thus, the convergence rate of Newton’s method is of order at
least two.
For a quadratic function, as in equation (4.1), steepest descent takes many steps to converge, while Newton’s
method reaches the minimum in one step. However, this single iteration in Newton’s method is more expensive,
because it requires both the gradient g
k
and the Hessian Q
k
to be evaluated, for a total of n +
_
n
2
_
derivatives. In
addition, the Hessian must be inverted, or, at least, system (4.15) must be solved. For very large problems, in which
the dimension n of x is thousands or more, storing and manipulating a Hessian can be prohibitive. In contrast, steepest
descent requires the gradient g
k
for selecting the step direction p
k
, and a line search in the direction p
k
to find the
step size. The method of conjugate gradients, discussed in the next section, is motivated by the desire to accelerate
convergence with respect to the steepest descent method, but without paying the storage cost of Newton’s method.
4.3 Conjugate Gradients
Newton’s method converges faster (quadratically) than steepest descent (linear convergence rate) because it uses more
information about the function f being minimized. Steepest descent locally approximates the function with planes,
because it only uses gradient information. All it can do is to go downhill. Newton’s method approximates f with
paraboloids, and then jumps at every iteration to the lowest point of the current approximation. The bottom line is that
fast convergence requires work that is equivalent to evaluating the Hessian of f.
Prima facie, the method of conjugate gradients discussed in this section seems to violate this principle: it achieves
fast, superlinear convergence, similarly to Newton’s method, but it only requires gradient information. This paradox,
however, is only apparent. Conjugate gradients works by taking n steps for each of the steps in Newton’s method.
It effectively solves the linear system (4.2) of Newton’s method, but it does so by a sequence of n one-dimensional
minimizations, each requiring one gradient computation and one line search.
Overall, the work done by conjugate gradients is equivalent to that done by Newton’s method. However, system
(4.2) is never constructed explicitly, and the matrix Q is never stored. This is very important in cases where x has
thousands or even millions of components. These high-dimensional problems arise typically from the discretization
of partial differential equations. Say for instance that we want to compute the motion of points in an image as a
consequence of camera motion. Partial differential equations relate image intensities over space and time to the motion
of the underlying image features. At every pixel in the image, this motion, called the motion field, is represented by
a vector whose magnitude and direction describe the velocity of the image feature at that pixel. Thus, if an image
has, say, a quarter of a million pixels, there are n = 500, 000 unknown motion field values. Storing and inverting a
500, 000 ×500, 000 Hessian is out of the question. In cases like these, conjugate gradients saves the day.
The conjugate gradients method described in these notes is the so-called Polak-Ribi` ere variation. It will be intro-
duced in three steps. First, it will be developed for the simple case of minimizing a quadratic function with positive-
definite and known Hessian. This quadratic function f(x) was introduced in equation (4.1). We know that in this case
minimizing f(x) is equivalent to solving the linear system (4.2). Rather than an iterative method, conjugate gradients
is a direct method for the quadratic case. This means that the number of iterations is fixed. Specifically, the method
converges to the solution in n steps, where n is the number of components of x. Because of the equivalence with
a linear system, conjugate gradients for the quadratic case can also be seen as an alternative method for solving a
linear system, although the version presented here will only work if the matrix of the system is symmetric and positive
definite.
4.3. CONJUGATE GRADIENTS 49
Second, the assumption that the Hessian Q in expression (4.1) is known will be removed. As discussed above, this
is the main reason for using conjugate gradients.
Third, the conjugate gradients method will be extended to general functions f(x). In this case, the method is no
longer direct, but iterative, and the cost of finding the minimum depends on the desired accuracy. This occurs because
the Hessian of f is no longer a constant, as it was in the quadratic case. As a consequence, a certain property that holds
in the quadratic case is now valid only approximately. In spite of this, the convergence rate of conjugate gradients is
superlinear, somewhere between Newton’s method and steepest descent. Finding tight bounds for the convergence rate
of conjugate gradients is hard, and we will omit this proof. We rely instead on the intuition that conjugate gradients
solves system (4.2), and that the quadratic approximation becomes more and more valid as the algorithm converges
to the minimum. If the function f starts to behave like a quadratic function early, that is, if f is nearly quadratic in a
large neighborhood of the minimum, convergence is fast, as it requires close to the n steps that are necessary in the
quadratic case, and each of the steps is simple. This combination of fast convergence, modest storage requirements,
and low computational cost per iteration explains the popularity of conjugate gradients methods for the optimization
of functions of a large number of variables.
4.3.1 The Quadratic Case
Suppose that we want to minimize the quadratic function
f(x) = c + a
T
x +
1
2
x
T
Qx (4.16)
where Q is a symmetric, positive definite matrix, and x has n components. As we saw in our discussion of steepest
descent, the minimum x

is the solution to the linear system
Qx = −a . (4.17)
We know how to solve such a system. However, all the methods we have seen so far involve explicit manipulation
of the matrix Q. We now consider an alternative solution method that does not need Q, but only the quantity
g
k
= Qx
k
+ a
that is, the gradient of f(x), evaluated at n different points x
1
, . . . , x
n
. We will see that the conjugate gradients method
requires n gradient evaluations and n line searches in lieu of each n ×n matrix inversion in Newton’s method.
Formal proofs can be found in Elijah Polak, Optimization — Algorithms and consistent approximations, Springer,
NY, 1997. The arguments offered below appeal to intuition.
Consider the case n = 3, in which the variable x in f(x) is a three-dimensional vector. Then the quadratic function
f(x) is constant over ellipsoids, called isosurfaces, centered at the minimum x

. How can we start from a point x
0
on one of these ellipsoids and reach x

by a finite sequence of one-dimensional searches? In connection with steepest
descent, we noticed that for poorly conditioned Hessians orthogonal directions lead to many small steps, that is, to
slow convergence.
When the ellipsoids are spheres, on the other hand, this works much better. The first step takes from x
0
to x
1
, and
the line between x
0
and x
1
is tangent to an isosurface at x
1
. The next step is in the direction of the gradient, so that
the new direction p
1
is orthogonal to the previous direction p
0
. This would then take us to x

right away. Suppose
however that we cannot afford to compute this special direction p
1
orthogonal to p
0
, but that we can only compute
some direction p
1
orthogonal to p
0
(there is an n − 1-dimensional space of such directions!). It is easy to see that in
that case n steps will take us to x

. In fact, since isosurfaces are spheres, each line minimization is independent of the
others: The first step yields the minimum in the space spanned by p
0
, the second step then yields the minimum in the
space spanned by p
0
and p
1
, and so forth. After n steps we must be done, since p
0
. . . , p
n−1
span the whole space.
In summary, any set of orthogonal directions, with a line search in each direction, will lead to the minimum for
spherical isosurfaces. Given an arbitrary set of ellipsoidal isosurfaces, there is a one-to-one mapping with a spherical
system: if Q = UΣU
T
is the SVD of the symmetric, positive definite matrix Q, then we can write
1
2
x
T
Qx =
1
2
y
T
y
50 CHAPTER 4. FUNCTION OPTIMIZATION
where
y = Σ
1/2
U
T
x . (4.18)
Consequently, there must be a condition for the original problem (in terms of Q) that is equivalent to orthogonality for
the spherical problem. If two directions q
i
and q
j
are orthogonal in the spherical context, that is, if
q
T
i
q
j
= 0 ,
what does this translate into in terms of the directions p
i
and p
j
for the ellipsoidal problem? We have
q
i,j
= Σ
1/2
U
T
p
i,j
,
so that orthogonality for q
i,j
becomes
p
T
i

1/2
Σ
1/2
U
T
p
j
= 0
or
p
T
i
Qp
j
= 0 . (4.19)
This condition is called Q-conjugacy, or Q-orthogonality: if equation (4.19) holds, then p
i
and p
j
are said to be
Q-conjugate or Q-orthogonal to each other. We will henceforth simply say “conjugate” for brevity.
In summary, if we can find n directions p
0
, . . . , p
n−1
that are mutually conjugate, and if we do line minimization
along each direction p
k
, we reach the minimum in at most n steps. Of course, we cannot use the transformation (4.18)
in the algorithm, because Σ and especially U
T
are too large. So now we need to find a method for generating n
conjugate directions without using either Q or its SVD. We do this in two steps. First, we find conjugate directions
whose definitions do involve Q. Then, in the next subsection, we rewrite these expressions without Q.
Here is the procedure, due to Hestenes and Stiefel (Methods of conjugate gradients for solving linear systems, J.
Res. Bureau National Standards, section B, Vol 49, pp. 409-436, 1952), which also incorporates the steps from x
0
to
x
n
:
g
0
= g(x
0
)
p
0
= −g
0
for k = 0 . . . , n −1
α
k
= arg min
α≥0
f(x
k
+αp
k
)
x
k+1
= x
k

k
p
k
g
k+1
= g(x
k+1
)
γ
k
=
g
T
k+1
Qp
k
p
T
k
Qp
k
p
k+1
= −g
k+1

k
p
k
end
where
g
k
= g(x
k
) =
∂f
∂x
¸
¸
¸
¸
x=x
k
is the gradient of f at x
k
.
It is simple to see that p
k
and p
k+1
are conjugate. In fact,
p
T
k
Qp
k+1
= p
T
k
Q(−g
k+1

k
p
k
)
= −p
T
k
Qg
k+1
+
g
T
k+1
Qp
k
p
T
k
Qp
k
p
T
k
Qp
k
= −p
T
k
Qg
k+1
+ g
T
k+1
Qp
k
= 0 .
It is somewhat more cumbersome to show that p
i
and p
k+1
for i = 0, . . . , k are also conjugate. This can be done by
induction. The proof is based on the observation that the vectors p
k
are found by a generalization of Gram-Schmidt
(theorem 2.4.2) to produce conjugate rather than orthogonal vectors. Details can be found in Polak’s book mentioned
earlier.
4.3. CONJUGATE GRADIENTS 51
4.3.2 Removing the Hessian
The algorithm shown in the previous subsection is a correct conjugate gradients algorithm. However, it is computa-
tionally inadequate because the expression for γ
k
contains the Hessian Q, which is too large. We now show that γ
k
can be rewritten in terms of the gradient values g
k
and g
k+1
only. To this end, we notice that
g
k+1
= g
k

k
Qp
k
,
or
α
k
Qp
k
= g
k+1
−g
k
.
In fact,
g(x) = a +Qx
so that
g
k+1
= g(x
k+1
) = g(x
k

k
p
k
) = a +Q(x
k

k
p
k
) = g
k

k
Qp
k
.
We can therefore write
γ
k
=
g
T
k+1
Qp
k
p
T
k
Qp
k
=
g
T
k+1
α
k
Qp
k
p
T
k
α
k
Qp
k
=
g
T
k+1
(g
k+1
−g
k
)
p
T
k
(g
k+1
−g
k
)
,
and Q has disappeared.
This expression for γ
k
can be further simplified by noticing that
p
T
k
g
k+1
= 0
because the line along p
k
is tangent to an isosurface at x
k+1
, while the gradient g
k+1
is orthogonal to the isosurface at
x
k+1
. Similarly,
p
T
k−1
g
k
= 0 .
Then, the denominator of γ
k
becomes
p
T
k
(g
k+1
−g
k
) = −p
T
k
g
k
= (g
k
−γ
k−1
p
k−1
)
T
g
k
= g
T
k
g
k
.
In conclusion, we obtain the Polak-Ribi` ere formula
γ
k
=
g
T
k+1
(g
k+1
−g
k
)
g
T
k
g
k
.
4.3.3 Extension to General Functions
We now know how to minimize the quadratic function (4.16) in n steps, without ever constructing the Hessian explic-
itly. When the function f(x) is arbitrary, the same algorithm can be used.
However, n iterations will not suffice. In fact, the Hessian, which was constant for the quadratic case, now is a
function of x
k
. Strictly speaking, we then lose conjugacy, since p
k
and p
k+1
are associated to different Hessians.
However, as the algorithm approaches the minimum x

, the quadratic approximation becomes more and more valid,
and a few cycles of n iterations each will achieve convergence.
52 CHAPTER 4. FUNCTION OPTIMIZATION
Chapter 5
Eigenvalues and Eigenvectors
Given a linear transformation
b = Ax ,
the singular value decomposition A = UΣV
T
of A transforms the domain of the transformation via the matrix V
T
and its range via the matrix U
T
so that the transformed system is diagonal. In fact, the equation b = UΣV
T
x can be
written as follows
U
T
b = ΣV
T
x ,
that is,
c = Σy
where
y = V
T
x and c = U
T
b ,
and where Σ is diagonal. This is a fundamental transformation to use whenever the domain and the range of A are
separate spaces. Often, however, domain and range are intimately related to one another even independently of the
transformation A. The most important example is perhaps that of a system of linear differential equations, of the form
˙ x = Ax
where Ais n×n. For this equation, the fact that Ais square is not a coincidence. In fact, x is assumed to be a function
of some real scalar variable t (often time), and ˙ x is the derivative of x with respect to t:
˙ x =
dx
dt
.
In other words, there is an intimate, pre-existing relation between x and ˙ x, and one cannot change coordinates for x
without also changing those for ˙ x accordingly. In fact, if V is an orthogonal matrix and we define
y = V
T
x ,
then the definition of ˙ x forces us to transform ˙ x by V
T
as well:
d V
T
x
dt
= V
T
dx
dt
= V
T
˙ x .
In brief, the SVD does nothing useful for systems of linear differential equations, because it diagonalizes Aby two
different transformations, one for the domain and one for the range, while we need a single transformation. Ideally,
we would like to find an orthogonal matrix S and a diagonal matrix Λ such that
A = SΛS
T
(5.1)
53
54 CHAPTER 5. EIGENVALUES AND EIGENVECTORS
so that if we define
y = S
T
x
we can write the equivalent but diagonal differential system
˙ y = Λy .
This is now much easier to handle, because it is a system of n independent, scalar differential equations, which can be
solved separately. The solutions can then be recombined through
x = Sy .
We will see all of this in greater detail soon.
Unfortunately, writing A in the form (5.1) is not always possible. This stands to reason, because now we are
imposing stronger constraints on the terms of the decomposition. It is like doing an SVD but with the additional
constraint U = V . If we refer back to figure 3.1, now the circle and the ellipse live in the same space, and the
constraint U = V implies that the vectors v
i
on the circle that map into the axes σ
i
u
i
of the ellipse are parallel to the
axes themselves. This will only occur for very special matrices.
In order to make a decomposition like (5.1) possible, we weaken the constraints in several ways:
• the elements of S and Λ are allowed to be complex, rather than real;
• the elements on the diagonal of Λ are allowed to be negative; in fact, they can be even non-real;
• S is required to be only invertible, rather than orthogonal.
To distinguish invertible from orthogonal matrices we use the symbol Q for invertible and S for orthogonal. In some
cases, it will be possible to diagonalize A by orthogonal transformations S and S
T
. Finally, for complex matrices we
generalize the notion of transpose by introducing the Hermitian operator: The matrix Q
H
(pronounced “QHermitian”)
is defined to be the complex conjugate of the transpose of Q. If Q happens to be real, conjugate transposition becomes
simply transposition, so the Hermitian is a generalization of the transpose. A matrix S is said to be unitary if
S
H
S = SS
H
= I ,
so unitary generalizes orthogonal for complex matrices. Unitary matrices merely rotate or flip vectors, in the sense
that they do not alter the vectors’ norms. For complex vectors, the norm squared is defined as
x
2
= x
H
x ,
and if S is unitary we have
Sx
2
= x
H
S
H
Sx = x
H
x = x
2
.
Furthermore, if x
1
and x
2
are mutually orthogonal, in the sense that
x
H
1
x
2
= 0 ,
then Sx
1
and Sx
2
are orthogonal as well:
x
H
1
S
H
Sx
2
= x
H
1
x
2
= 0 .
In contrast, a nonunitary transformation Q can change the norms of vectors, as well as the inner products between
vectors. A matrix that is equal to its Hermitian is called a Hermitian matrix.
In summary, in order to diagonalize a square matrix A from a system of linear differential equations we generally
look for a decomposition of A of the form
A = QΛQ
−1
(5.2)
55
where Q and Λ are complex, Q is invertible, and Λ is diagonal. For some special matrices, this may specialize to
A = SΛS
H
with unitary S.
Whenever two matrices A and B, diagonal or not, are related by
A = QBQ
−1
,
they are said to be similar to each other, and the transformation of B into A (and vice versa) is called a similarity
transformation.
The equation A = QΛQ
−1
can be rewritten as follows:
AQ = QΛ
or separately for every column of Q as follows:
Aq
i
= λ
i
q
i
(5.3)
where
Q =
_
q
1
· · · q
n
¸
and Λ = diag(λ
1
, . . . , λ
n
) .
Thus, the columns of q
i
of Q and the diagonal entries λ
i
of Λ are solutions of the eigenvalue/eigenvector equation
Ax = λx , (5.4)
which is how eigenvalues and eigenvectors are usually introduced. In contrast, we have derived this equation from the
requirement of diagonalizing a matrix by a similarity transformation. The columns of Q are called eigenvectors, and
the diagonal entries of Λ are called eigenvalues.
−2 −1 0 1 2
−2
−1.5
−1
−0.5
0
0.5
1
1.5
2
Figure 5.1: Effect of the transformation (5.5) on a sample of points on the unit circle. The dashed lines are vectors that
do not change direction under the transformation.
That real eigenvectors and eigenvalues do not always exist can be clarified by considering the eigenvalue problem
from a geometrical point of view in the n = 2 case. As we know, an invertible linear transformation transforms the
unit circle into an ellipse. Each point on the unit circle is transformed into some point on the ellipse. Figure 5.1 shows
the effect of the transformation represented by the matrix
A =
_
2/3 4/

3
0 2
_
(5.5)
56 CHAPTER 5. EIGENVALUES AND EIGENVECTORS
for a sample of points on the unit circle. Notice that there are many transformations that map the unit circle into the
same ellipse. In fact, the circle in figure 5.1 can be rotated, pulling the solid lines along. Each rotation yields another
matrix A, but the resulting ellipse is unchanged. In other words, the curve-to-curve transformation from circle to
ellipse is unique, but the point-to-point transformation is not. Matrices represent point-to-point transformations.
The eigenvalue problem amounts to finding axes q
1
, q
2
that are mapped into themselves by the original transfor-
mation A (see equation (5.3)). In figure 5.1, the two eigenvectors are shown as dashed lines. Notice that they do not
correspond to the axes of the ellipse, and that they are not orthogonal. Equation (5.4) is homogeneous in x, so x can
be assumed to be a unit vector without loss of generality.
Given that the directions of the input vectors are generally changed by the transformation A, as evident from figure
5.1, it is not obvious whether the eigenvalue problem admits a solution at all. We will see that the answer depends
on the matrix A, and that a rather diverse array of situations may arise. In some cases, the eigenvalues and their
eigenvectors exist, but they are complex. The geometric intuition is hidden, and the problem is best treated as an
algebraic one. In other cases, all eigenvalues exist, perhaps all real, but not enough eigenvectors can be found, and the
matrix A cannot be diagonalized. In particularly good cases, there are n real, orthonormal eigenvectors. In bad cases,
we have to give up the idea of diagonalizing A, and we can only triangularize it. This turns out to be good enough for
solving linear differential systems, just as triangularization was sufficient for solving linear algebraic systems.
5.1 Computing Eigenvalues and Eigenvectors Algebraically
Let us rewrite the eigenvalue equation
Ax = λx
as follows:
(A−λI)x = 0 . (5.6)
This is a homogeneous, square system of equations, which admits nontrivial solutions iff the matrix A − λI is rank-
deficient. A square matrix B is rank-deficient iff its determinant,
det(B) =
_
b
11
if B is 1 ×1

n
i=1
(−1)
i+1
b
i1
det(B
i1
) otherwise
is zero. In this expression, B
ij
is the algebraic complement of entry b
ij
, defined as the (n − 1) × (n − 1) matrix
obtained by removing row i and column j from B.
Volumes have been written about the properties of the determinant. For our purposes, it is sufficient to recall the
following properties from linear algebra:
• det(B) = det(B
T
);
• det(
_
b
1
· · · b
n
¸
) = 0 iff b
1
, . . . , b
n
are linearly dependent;
• det(
_
b
1
· · · b
i
· · · b
j
· · · b
n
¸
) = −det(
_
b
1
· · · b
j
· · · b
i
· · · b
n
¸
);
• det(BC) = det(B) det(C).
Thus, for system (5.6) to admit nontrivial solutions, we need
det(A−λI) = 0 . (5.7)
From the definition of determinant, it follows, by very simple induction, that the left-hand side of equation (5.7)
is a polynomial of degree n in λ, and that the coefficient of λ
n
is 1. Therefore, equation (5.7), which is called the
characteristic equation of A, has n complex solutions, in the sense that
det(A−λI) = (−1)
n
(λ −λ
1
) · . . . · (λ −λ
n
)
where some of the λ
i
may coincide. In other words, an n × n matrix has at most n distinct eigenvalues. The case of
exactly n distinct eigenvalues is of particular interest, because of the following results.
5.1. COMPUTING EIGENVALUES AND EIGENVECTORS ALGEBRAICALLY 57
Theorem 5.1.1 Eigenvectors x
1
, . . . , x
k
corresponding to distinct eigenvalues λ
1
, . . . , λ
k
are linearly independent.
Proof. Suppose that c
1
x
1
+. . . +c
k
x
k
= 0 where the x
i
are eigenvectors of a matrix A. We need to show that then
c
1
= . . . = c
k
= 0. By multiplying by A we obtain
c
1
Ax
1
+. . . +c
k
Ax
k
= 0
and because x
1
, . . . , x
k
are eigenvectors corresponding to eigenvalues λ
1
, . . . , λ
k
, we have
c
1
λ
1
x
1
+. . . +c
k
λ
k
x
k
= 0 . (5.8)
However, from
c
1
x
1
+. . . +c
k
x
k
= 0
we also have
c
1
λ
k
x
1
+. . . +c
k
λ
k
x
k
= 0
and subtracting this equation from equation (5.8) we have
c
1

1
−λ
k
)x
1
+. . . +c
k−1

k−1
−λ
k
)x
k−1
= 0 .
Thus, we have reduced the summation to one containing k − 1 terms. Since all λ
i
are distinct, the differences in
parentheses are all nonzero, and we can replace each x
i
by x

i
= (λ
i
−λ
k
)x
i
, which is still an eigenvector of A:
c
1
x

1
+. . . +c
k−1
x

k−1
= 0 .
We can repeat this procedure until only one term remains, and this forces c
1
= 0, so that
c
2
x
2
+. . . +c
k
x
k
= 0
This entire argument can be repeated for the last equation, therefore forcing c
2
= 0, and so forth.
In summary, the equation c
1
x
1
+. . . +c
k
x
k
= 0 implies that c
1
= . . . = c
k
= 0, that is, that the vectors x
1
, . . . , x
k
are linearly independent. ∆
For Hermitian matrices (and therefore for real symmetric matrices as well), the situation is even better.
Theorem 5.1.2 A Hermitian matrix has real eigenvalues.
Proof. A matrix A is Hermitian iff A = A
H
. Let λ and x be an eigenvalue of A and a corresponding eigenvector:
Ax = λx . (5.9)
By taking the Hermitian we obtain
x
H
A
H
= λ

x
H
.
Since A = A
H
, the last equation can be rewritten as follows:
x
H
A = λ

x
H
. (5.10)
If we multiply equation (5.9) from the left by x
H
and equation (5.10) from the right by x, we obtain
x
H
Ax = λx
H
x
x
H
Ax = λ

x
H
x
which implies that
λx
H
x = λ

x
H
x
58 CHAPTER 5. EIGENVALUES AND EIGENVECTORS
Since x is an eigenvector, the scalar x
H
x is nonzero, so that we have
λ = λ

as promised. ∆
Corollary 5.1.3 A real and symmetric matrix has real eigenvalues.
Proof. A real and symmetric matrix is Hermitian. ∆
Theorem 5.1.4 Eigenvectors corresponding to distinct eigenvalues of a Hermitian matrix are mutually orthogonal.
Proof. Let λ and µ be two distinct eigenvalues of A, and let x and y be corresponding eigenvectors:
Ax = λx
Ay = µy ⇒ y
H
A = µy
H
because A = A
H
and from theorem 5.1.2 µ = µ

. If we multiply these two equations by y
H
from the left and x from
the right, respectively, we obtain
y
H
Ax = λy
H
x
y
H
Ax = µy
H
x ,
which implies
λy
H
x = µy
H
x
or
(λ −µ)y
H
x = 0 .
Since the two eigenvalues are distinct, λ −µ is nonzero, and we must have y
H
x = 0. ∆
Corollary 5.1.5 An n ×n Hermitian matrix with n distinct eigenvalues admits n orthonormal eigenvectors.
Proof. From theorem 5.1.4, the eigenvectors of an n×n Hermitian matrix with n distinct eigenvalues are all mutu-
ally orthogonal. Since the eigenvalue equation Ax = λx is homogeneous in x, the vector x can be normalized without
violating the equation. Consequently, the eigenvectors can be made to be orthonormal. ∆
In summary, any square matrix with n distinct eigenvalues can be diagonalized by a similarity transformation, and
any square Hermitian matrix with n distinct eigenvalues can be diagonalized by a unitary similarity transformation.
Notice that the converse is not true: a matrix can have coincident eigenvalues and still admit n independent, and
even orthonormal, eigenvectors. For instance, the n × n identity matrix has n equal eigenvalues but n orthonormal
eigenvectors (which can be chosen in infinitely many ways).
The examples in section 5.2 show that when some eigenvalues coincide, rather diverse situations can arise concern-
ing the eigenvectors. First, however, we point out a simple but fundamental fact about the eigenvalues of a triangular
matrix.
Theorem 5.1.6 The determinant of a triangular matrix is the product of the elements on its diagonal.
5.2. GOOD AND BAD MATRICES 59
Proof. This follows immediately from the definition of determinant. Without loss of generality, we can assume a
triangular matrix B to be upper-triangular, for otherwise we can repeat the argument for the transpose, which because
of the properties above has the same eigenvalues. Then, the only possibly nonzero b
i1
of the matrix B is b
11
, and the
summation in the definition of determinant given above reduces to a single term:
det(B) =
_
b
11
if B is 1 ×1
b
11
det(B
11
) otherwise
.
By repeating the argument for B
11
and so forth until we are left with a single scalar, we obtain
det(B) = b
11
· . . . · b
nn
.

Corollary 5.1.7 The eigenvalues of a triangular matrix are the elements on its diagonal.
Proof. The eigenvalues of a matrix A are the solutions of the equation
det(A−λI) = 0 .
If A is triangular, so is B = A−λI, and from the previous theorem we obtain
det(A−λI) = (a
11
−λ) · . . . · (a
nn
−λ)
which is equal to zero for
λ = a
11
, . . . , a
nn
.

Note that diagonal matrices are triangular, so this result holds for diagonal matrices as well.
5.2 Good and Bad Matrices
Solving differential equations becomes much easier when matrices have a full set of orthonormal eigenvectors. For
instance, the matrix
A =
_
2 0
0 1
_
(5.11)
has eigenvalues 2 and 1 and eigenvectors
s
1
=
_
1
0
_
s
2
=
_
0
1
_
.
Matrices with n orthonormal eigenvectors are called normal. Normal matrices are good news, because then the
n ×n system of differential equations
˙ x = Ax
has solution
x(t) =
n

i=1
c
i
s
i
e
λit
= S
_
¸
_
e
λ
1
t
.
.
.
e
λnt
_
¸
_c
60 CHAPTER 5. EIGENVALUES AND EIGENVECTORS
where S = [s
1
· · · s
n
] are the eigenvectors, λ
i
are the eigenvalues, and the vector c of constants c
i
is
c = S
H
x(0) .
More compactly,
x(t) = S
_
¸
_
e
λ1t
.
.
.
e
λ
n
t
_
¸
_S
H
x(0) .
Fortunately these matrices occur frequently in practice. However, not all matrices are as good as these. First, there
may still be a complete set of n eigenvectors, but they may not be orthonormal. An example of such a matrix is
_
2 −1
0 1
_
which has eigenvalues 2 and 1 and nonorthogonal eigenvectors
q
1
=
_
1
0
_
q
2
=

2
2
_
1
1
_
.
This is conceptually only a slight problem, because the unitary matrix S is replaced by an invertible matrix Q, and the
solution becomes
x(t) = Q
_
¸
_
e
λ1t
.
.
.
e
λnt
_
¸
_Q
−1
x(0) .
Computationally this is more expensive, because a computation of a Hermitian is replaced by a matrix inversion.
However, things can be worse yet, and a full set of eigenvectors may fail to exist, as we now show.
A necessary condition for an n ×n matrix to be defective, that is, to have fewer than n eigenvectors, is that it have
repeated eigenvalues. In fact, we have seen (theorem 5.1.1) that a matrix with distinct eigenvalues (zero or nonzero
does not matter) has a full set of eigenvectors (perhaps nonorthogonal, but independent). The simplest example of a
defective matrix is
_
0 1
0 0
_
which has double eigenvalue 0 and only eigenvector [1 0]
T
, while
_
3 1
0 3
_
has double eigenvalue 3 and only eigenvector [1 0]
T
, so zero eigenvalues are not the problem.
However, repeated eigenvalues are not a sufficient condition for defectiveness, as the identity matrix proves.
How bad can a matrix be? Here is a matrix that is singular, has fewer than n eigenvectors, and the eigenvectors it
has are not orthogonal. It belongs to the scum of all matrices:
A =
_
_
0 2 −1
0 2 1
0 0 2
_
_
.
Its eigenvalues are 0, because the matrix is singular, and 2, repeated twice. A has to have a repeated eigenvalue if it is
to be defective. Its two eigenvectors are
q
1
=
_
_
1
0
0
_
_
, q
2
=

2
2
_
_
1
1
0
_
_
corresponding to eigenvalues 0 and 2 in this order, and there is no q
3
. Furthermore, q
1
and q
2
are not orthogonal to
each other.
5.3. COMPUTING EIGENVALUES AND EIGENVECTORS NUMERICALLY 61
5.3 Computing Eigenvalues and Eigenvectors Numerically
The examples above have shown that not every n × n matrix admits n independent eigenvectors, so some matrices
cannot be diagonalized by similarity transformations. Fortunately, these matrices can be triangularized by similarity
transformations, as we now show. We will show later on that this allows solving systems of linear differential equations
regardless of the structure of the system’s matrix of coefficients.
It is important to notice that if a matrix A is triangularized by similarity transformations,
T = Q
−1
AQ ,
then the eigenvalues of the triangular matrix T are equal to those of the original matrix A. In fact, if
Ax = λx ,
then
QTQ
−1
x = λx ,
that is,
Ty = λy
where
y = Q
−1
x ,
so λ is also an eigenvalue for T. The eigenvectors, however, are changed according to the last equation.
The Schur decomposition does even better, since it triangularizes any square matrix A by a unitary (possibly
complex) transformation:
T = S
H
AS .
This transformation is equivalent to factoring A into the product
A = STS
H
,
and this product is called the Schur decomposition of A. Numerically stable and efficient algorithms exist for the Schur
decomposition. In this note, we will not study these algorithms, but only show that all square matrices admit a Schur
decomposition.
5.3.1 Rotations into the x
1
Axis
An important preliminary fact concerns vector rotations. Let e
1
be the first column of the identity matrix. It is
intuitively obvious that any nonzero real vector x can be rotated into a vector parallel to e
1
. Formally, take any
orthogonal matrix S whose first column is
s
1
=
x
x
.
Since s
T
1
x = x
T
x/x = x, and since all the other s
j
are orthogonal to s
1
, we have
S
T
x =
_
¸
_
s
T
1
.
.
.
s
T
n
_
¸
_x =
_
¸
_
s
T
1
x
.
.
.
s
T
n
x
_
¸
_ =
_
¸
¸
¸
_
x
0
.
.
.
0
_
¸
¸
¸
_
which is parallel to e
1
as desired. It may be less obvious that a complex vector x can be transformed into a real vector
parallel to e
1
by a unitary transformation. But the trick is the same: let
s
1
=
x
x
.
62 CHAPTER 5. EIGENVALUES AND EIGENVECTORS
Now s
1
may be complex. We have s
H
1
x = x
H
x/x = x, and
S
H
x =
_
¸
_
s
H
1
.
.
.
s
H
n
_
¸
_x =
_
¸
_
s
H
1
x
.
.
.
s
H
n
x
_
¸
_ =
_
¸
¸
¸
_
x
0
.
.
.
0
_
¸
¸
¸
_
,
just about like before. We are now ready to triangularize an arbitrary square matrix A.
5.3.2 The Schur Decomposition
The Schur decomposition theorem is the cornerstone of eigenvalue computations. It states that any square matrix can
be triangularized by unitary transformations. The diagonal elements of a triangular matrix are its eigenvalues, and
unitary transformations preserve eigenvalues. Consequently, if the Schur decomposition of a matrix can be computed,
its eigenvalues can be determined. Moreover, as we will see later, a system of linear differential equations can be
solved regardless of the structure of the matrix of its coefficients.
Lemma 5.3.1 If A is an n ×n matrix and λ and x are an eigenvalue of A and its corresponding eigenvector,
Ax = λx (5.12)
then there is a transformation
T = U
H
AU
where U is a unitary, n ×n matrix, such that
T =
_
¸
¸
¸
_
λ
0
.
.
.
0
C
_
¸
¸
¸
_
.
Proof. Let U be a unitary transformation that transforms the (possibly complex) eigenvector x of A into a real
vector on the x
1
axis:
x = U
_
¸
¸
¸
_
r
0
.
.
.
0
_
¸
¸
¸
_
where r is the nonzero norm of x. By substituting this into (5.12) and rearranging we have
AU
_
¸
¸
¸
_
r
0
.
.
.
0
_
¸
¸
¸
_
= λU
_
¸
¸
¸
_
r
0
.
.
.
0
_
¸
¸
¸
_
U
H
AU
_
¸
¸
¸
_
r
0
.
.
.
0
_
¸
¸
¸
_
= λ
_
¸
¸
¸
_
r
0
.
.
.
0
_
¸
¸
¸
_
U
H
AU
_
¸
¸
¸
_
1
0
.
.
.
0
_
¸
¸
¸
_
= λ
_
¸
¸
¸
_
1
0
.
.
.
0
_
¸
¸
¸
_
5.3. COMPUTING EIGENVALUES AND EIGENVECTORS NUMERICALLY 63
T
_
¸
¸
¸
_
1
0
.
.
.
0
_
¸
¸
¸
_
=
_
¸
¸
¸
_
λ
0
.
.
.
0
_
¸
¸
¸
_
.
The last left-hand side is the first column of T, and the corresponding right-hand side is of the form required by the
lemma. ∆
Theorem 5.3.2 (Schur) If A is any n ×n matrix then there exists a unitary n ×n matrix S such that
S
H
AS = T
where T is triangular. Furthermore, S can be chosen so that the eigenvalues λ
i
of A appear in any order along the
diagonal of T.
Proof. By induction. The theorem obviously holds for n = 1:
1 A1 = A .
Suppose it holds for all matrices of order n −1. Then from the lemma there exists a unitary U such that
U
H
AU =
_
¸
¸
¸
_
λ
0
.
.
.
0
C
_
¸
¸
¸
_
where λ is any eigenvalue of A. Partition C into a row vector and an (n −1) ×(n −1) matrix G:
C =
_
w
H
G
_
.
By the inductive hypothesis, there is a unitary matrix V such that V
H
GV is a Schur decomposition of G. Let
S = U
_
¸
¸
¸
_
1 0 · · · 0
0
.
.
.
0
V
_
¸
¸
¸
_
.
Clearly, S is a unitary matrix, and S
H
AS is upper-triangular. Since the elements on the diagonal of a triangular matrix
are the eigenvalues, S
H
AS is the Schur decomposition of A. Because we can pick any eigenvalue as λ, the order of
eigenvalues can be chosen arbitrarily. ∆
This theorem does not say how to compute the Schur decomposition, only that it exists. Fortunately, there is a
stable and efficient algorithm to compute the Schur decomposition. This is the preferred way to compute eigenvalues
numerically.
64 CHAPTER 5. EIGENVALUES AND EIGENVECTORS
5.4 Eigenvalues/Vectors and Singular Values/Vectors
In this section we prove a few additional important properties of eigenvalues and eigenvectors. In the process, we also
establish a link between singular values/vectors and eigenvalues/vectors. While this link is very important, it is useful
to remember that eigenvalues/vectors and singular values/vectors are conceptually and factually very distinct entities
(recall figure 5.1).
First, a general relation between determinant and eigenvalues.
Theorem 5.4.1 The determinant of a matrix is equal to the product of its eigenvalues.
Proof. The proof is very simple, given the Schur decomposition. In fact, we know that the eigenvalues of a matrix
A are equal to those of the triangular matrix in the Schur decomposition of A. Furthermore, we know from theorem
5.1.6 that the determinant of a triangular matrix is the product of the elements on its diagonal. If we recall that a
unitary matrix has determinant 1 or -1, that the determinants of S and S
H
are the same, and that the determinant of a
product of matrices is equal to the product of the determinants, the proof is complete. ∆
We saw that an n ×n Hermitian matrix with n distinct eigenvalues admits n orthonormal eigenvectors (corollary
5.1.5). The assumption of distinct eigenvalues made the proof simple, but is otherwise unnecessary. In fact, now that
we have the Schur decomposition, we can state the following stronger result.
Theorem 5.4.2 (Spectral theorem) Every Hermitian matrix can be diagonalized by a unitary matrix, and every real
symmetric matrix can be diagonalized by an orthogonal matrix:
A = A
H
⇒ A = SΛS
H
A real, A = A
T
⇒ A = SΛS
T
, S real .
In either case, Λ is real and diagonal.
Proof. We already know that Hermitian matrices (and therefore real and symmetric ones) have real eigenvalues
(theorem 5.1.2), so Λ is real. Let now
A = STS
H
be the Schur decomposition of A. Since A is Hermitian, so is T. In fact, T = S
H
AS, and
T
H
= (S
H
AS)
H
= S
H
A
H
S = S
H
AS = T .
But the only way that T can be both triangular and Hermitian is for it to be diagonal, because 0

= 0. Thus, the
Schur decomposition of a Hermitian matrix is in fact a diagonalization, and this is the first equation of the theorem
(the diagonal of a Hermitian matrix must be real).
Let now Abe real and symmetric. All that is left to prove is that then its eigenvectors are real. But eigenvectors are
the solution of the homogeneous system (5.6), which is both real and rank-deficient, and therefore admits nontrivial
real solutions. Thus, S is real, and S
H
= S
T
. ∆
In other words, a Hermitian matrix, real or not, with distinct eigenvalues or not, has real eigenvalues and n or-
thonormal eigenvectors. If in addition the matrix is real, so are its eigenvectors.
We recall that a real matrix A such that for every nonzero x we have x
T
Ax > 0 is said to be positive definite. It is
positive semidefinite if for every nonzero x we have x
T
Ax ≥ 0. Notice that a positive definite matrix is also positive
semidefinite. Positive definite or semidefinite matrices arise in the solution of overconstrained linear systems, because
A
T
A is positive semidefinite for every A (lemma 5.4.5). They also occur in geometry through the equation of an
ellipsoid,
x
T
Qx = 1
5.4. EIGENVALUES/VECTORS AND SINGULAR VALUES/VECTORS 65
in which Q is positive definite. In physics, positive definite matrices are associated to quadratic forms x
T
Qx that
represent energies or second-order moments of mass or force distributions. Their physical meaning makes them
positive definite, or at least positive semidefinite (for instance, energies cannot be negative). The following result
relates eigenvalues/vectors with singular values/vectors for positive semidefinite matrices.
Theorem 5.4.3 The eigenvalues of a real, symmetric, positive semidefinite matrix A are equal to its singular values.
The eigenvectors of A are also its singular vectors, both left and right.
Proof. From the previous theorem, A = SΛS
T
, where both Λ and S are real. Furthermore, the entries in Λ are
nonnegative. In fact, from
As
i
= λs
i
we obtain
s
T
i
As
i
= s
T
i
λs
i
= λs
T
i
s
i
= λs
i

2
= λ .
If A is positive semidefinite, then x
T
Ax ≥ 0 for any nonzero x, and in particular s
T
i
As
i
≥ 0, so that λ ≥ 0.
But
A = SΛS
T
with nonnegative diagonal entries in Λ is the singular value decomposition A = UΣV
T
of A with Σ = Λ and
U = V = S. Recall that the eigenvalues in the Schur decomposition can be arranged in any desired order along the
diagonal. ∆
Theorem 5.4.4 A real, symmetric matrix is positive semidefinite iff all its eigenvalues are nonnegative. It is positive
definite iff all its eigenvalues are positive.
Proof. Theorem 5.4.3 implies one of the two directions: If A is real, symmetric, and positive semidefinite, then its
eigenvalues are nonnegative. If the proof of that theorem is repeated with the strict inequality, we also obtain that if A
is real, symmetric, and positive definite, then its eigenvalues are positive.
Conversely, we show that if all eigenvalues λ of a real and symmetric matrix A are positive (nonnegative) then A
is positive definite (semidefinite). To this end, let x be any nonzero vector. Since real and symmetric matrices have n
orthonormal eigenvectors (theorem 5.4.2), we can use these eigenvectors s
1
, . . . , s
n
as an orthonormal basis for R
n
,
and write
x = c
1
s
1
+. . . +c
n
s
n
with
c
i
= x
T
s
i
.
But then
x
T
Ax = x
T
A(c
1
s
1
+. . . +c
n
s
n
) = x
T
(c
1
As
1
+. . . +c
n
As
n
)
= x
T
(c
1
λ
1
s
1
+. . . +c
n
λ
n
s
n
) = c
1
λ
1
x
T
s
1
+. . . +c
n
λ
n
x
T
s
n
= λ
1
c
2
1
+. . . +λ
n
c
2
n
> 0 (or ≥ 0)
because the λ
i
are positive (nonnegative) and not all c
i
can be zero. Since x
T
Ax > 0 (or ≥ 0) for every nonzero x, A
is positive definite (semidefinite). ∆
Theorem 5.4.3 establishes one connection between eigenvalues/vectors and singular values/vectors: for symmetric,
positive definite matrices, the concepts coincide. This result can be used to introduce a less direct link, but for arbitrary
matrices.
Lemma 5.4.5 A
T
A is positive semidefinite.
66 CHAPTER 5. EIGENVALUES AND EIGENVECTORS
Proof. For any nonzero x we can write x
T
A
T
Ax = Ax
2
≥ 0. ∆
Theorem 5.4.6 The eigenvalues of A
T
A with m ≥ n are the squares of the singular values of A; the eigenvectors of
A
T
Aare the right singular vectors of A. Similarly, for m ≤ n, the eigenvalues of AA
T
are the squares of the singular
values of A, and the eigenvectors of AA
T
are the left singular vectors of A.
Proof. If m ≥ n and A = UΣV
T
is the SVD of A, we have
A
T
A = V ΣU
T
UΣV
T
= V Σ
2
V
T
which is in the required format to be a (diagonal) Schur decomposition with S = V and T = Λ = Σ
2
. Similarly, for
m ≤ n,
AA
T
= UΣV
T
V ΣU
T
= UΣ
2
U
T
is a Schur decomposition with S = U and T = Λ = Σ
2
. ∆
We have seen that important classes of matrices admit a full set of orthonormal eigenvectors. The theorem below
characterizes the class of all matrices with this property, that is, the class of all normal matrices. To prove the theorem,
we first need a lemma.
Lemma 5.4.7 If for an n ×n matrix B we have BB
H
= B
H
B, then for every i = 1, . . . , n, the norm of the i-th row
of B equals the norm of its i-th column.
Proof. From BB
H
= B
H
B we deduce
Bx
2
= x
H
B
H
Bx = x
H
BB
H
x = B
H
x
2
. (5.13)
If x = e
i
, the i-th column of the n × n identity matrix, Be
i
is the i-th column of B, and B
H
e
i
is the i-th column of
B
H
, which is the conjugate of the i-th row of B. Since conjugation does not change the norm of a vector, the equality
(5.13) implies that the i-th column of B has the same norm as the i-th row of B. ∆
Theorem 5.4.8 An n ×n matrix is normal if an only if it commutes with its Hermitian:
AA
H
= A
H
A .
Proof. Let A = STS
H
be the Schur decomposition of A. Then,
AA
H
= STS
H
ST
H
S
H
= STT
H
S
H
and A
H
A = ST
H
S
H
STS
H
= ST
H
TS
H
.
Because S is invertible (even unitary), we have AA
H
= A
H
A if and only if TT
H
= T
H
T.
However, a triangular matrix T for which TT
H
= T
H
T must be diagonal. In fact, from the lemma, the norm of
the i-th row of T is equal to the norm of its i-th column. Let i = 1. Then, the first column of T has norm |t
11
|. The
first row has first entry t
11
, so the only way that its norm can be |t
11
| is for all other entries in the first row to be zero.
We now proceed through i = 2, . . . , n, and reason similarly to conclude that T must be diagonal.
The converse is also obviously true: if T is diagonal, then TT
H
= T
H
T. Thus, AA
H
= A
H
A if and only if T is
diagonal, that is, if and only if A can be diagonalized by a unitary similarity transformation. This is the definition of a
normal matrix. ∆
5.4. EIGENVALUES/VECTORS AND SINGULAR VALUES/VECTORS 67
Corollary 5.4.9 A triangular, normal matrix must be diagonal.
Proof. We proved this in the proof of theorem 5.4.8. ∆
Checking that A
H
A = AA
H
is much easier than computing eigenvectors, so theorem 5.4.8 is a very useful
characterization of normal matrices. Notice that Hermitian (and therefore also real symmetric) matrices commute
trivially with their Hermitians, but so do, for instance, unitary (and therefore also real orthogonal) matrices:
UU
H
= U
H
U = I .
Thus, Hermitian, real symmetric, unitary, and orthogonal matrices are all normal.
68 CHAPTER 5. EIGENVALUES AND EIGENVECTORS
Chapter 6
Ordinary Differential Systems
In this chapter we use the theory developed in chapter 5 in order to solve systems of first-order linear differential
equations with constant coefficients. These systems have the following form:
˙ x = Ax + b(t) (6.1)
x(0) = x
0
(6.2)
where x = x(t) is an n-dimensional vector function of time t, the dot denotes differentiation, the coefficients a
ij
in
the n ×n matrix A are constant, and the vector function b(t) is a function of time. The equation (6.2), in which x
0
is
a known vector, defines the initial value of the solution.
First, we show that scalar differential equations of order greater than one can be reduced to systems of first-order
differential equations. Then, in section 6.2, we recall a general result for the solution of first-order differential systems
from the elementary theory of differential equations. In section 6.3, we make this result more specific by showing
that the solution to a homogeneous system is a linear combination of exponentials multiplied by polynomials in t.
This result is based on the Schur decomposition introduced in chapter 5, which is numerically preferable to the more
commonly used Jordan canonical form. Finally, in sections 6.4 and 6.5, we set up and solve a particular differential
system as an illustrative example.
6.1 Scalar Differential Equations of Order Higher than One
The first-order system (6.1) subsumes also the case of a scalar differential equation of order n, possibly greater than 1,
d
n
y
dt
n
+c
n−1
d
n−1
y
dt
n−1
+. . . +c
1
dy
dt
+c
0
y = b(t) . (6.3)
In fact, such an equation can be reduced to a first-order system of the form (6.1) by introducing the n-dimensional
vector
x =
_
¸
_
x
1
.
.
.
x
n
_
¸
_ =
_
¸
¸
¸
_
y
dy
dt
.
.
.
d
n−1
y
dt
n−1
_
¸
¸
¸
_
.
With this definition, we have
d
i
y
dt
i
= x
i+1
for i = 0, . . . , n −1
d
n
y
dt
n
=
dx
n
dt
,
69
70 CHAPTER 6. ORDINARY DIFFERENTIAL SYSTEMS
and x satisfies the additional n −1 equations
x
i+1
=
dx
i
dt
(6.4)
for i = 1, . . . , n − 1. If we write the original system (6.3) together with the n − 1 differential equations (6.4), we
obtain the first-order system
˙ x = Ax + b(t)
where
A =
_
¸
¸
¸
¸
¸
_
0 1 0 · · · 0
0 0 1 · · · 0
.
.
.
.
.
.
.
.
.
.
.
.
.
.
.
0 0 0 · · · 1
−c
0
−c
1
−c
2
· · · −c
n−1
_
¸
¸
¸
¸
¸
_
is the so-called companion matrix of (6.3) and
b(t) =
_
¸
¸
¸
¸
¸
_
0
0
.
.
.
0
b(t)
_
¸
¸
¸
¸
¸
_
.
6.2 General Solution of a Linear Differential System
We know from the general theory of differential equations that a general solution of system (6.1) with initial condition
(6.2) is given by
x(t) = x
h
(t) + x
p
(t)
where x
h
(t) is the solution of the homogeneous system
˙ x = Ax
x(0) = x
0
and x
p
(t) is a particular solution of
˙ x = Ax + b(t)
x(0) = 0 .
The two solution components x
h
and x
p
can be written by means of the matrix exponential, introduced in the following.
For the scalar exponential e
λt
we can write a Taylor series expansion
e
λt
= 1 +
λt
1!
+
λ
2
t
2
2!
+· · · =

j=0
λ
j
t
j
j!
.
Usually
1
, in calculus classes, the exponential is introduced by other means, and the Taylor series expansion above is
proven as a property.
For matrices, the exponential e
Z
of a matrix Z ∈ R
n×n
is instead defined by the infinite series expansion
e
Z
= I +
Z
1!
+
Z
2
2!
+· · · =

j=0
Z
j
j!
.
1
Not always. In some treatments, the exponential is defined through its Taylor series.
6.2. GENERAL SOLUTION OF A LINEAR DIFFERENTIAL SYSTEM 71
Here I is the n ×n identity matrix, and the general term Z
j
/j! is simply the matrix Z raised to the jth power divided
by the scalar j!. It turns out that this infinite sum converges (to an n×n matrix which we write as e
Z
) for every matrix
Z. Substituting Z = At gives
e
At
= I +
At
1!
+
A
2
t
2
2!
+
A
3
t
3
3!
+· · · =

j=0
A
j
t
j
j!
. (6.5)
Differentiating both sides of (6.5) gives
de
At
dt
= A+
A
2
t
1!
+
A
3
t
2
2!
+· · ·
= A
_
I +
At
1!
+
A
2
t
2
2!
+· · ·
_
de
At
dt
= Ae
At
.
Thus, for any vector w, the function x
h
(t) = e
At
w satisfies the homogeneous differential system
˙ x
h
= Ax
h
.
By using the initial values (6.2) we obtain w = x
0
, and
x
h
(t) = e
At
x(0) (6.6)
is a solution to the differential system (6.1) with b(t) = 0 and initial values (6.2). It can be shown that this solution is
unique.
From the elementary theory of differential equations, we also know that a particular solution to the nonhomoge-
neous (b(t) = 0) equation (6.1) is given by
x
p
(t) =
_
t
0
e
A(t−s)
b(s) ds .
This is easily verified, since by differentiating this expression for x
p
we obtain
˙ x
p
= Ae
At
_
t
0
e
−As
b(s) ds +e
At
e
−At
b(t) = Ax
p
+ b(t) ,
so x
p
satisfies equation (6.1).
In summary, we have the following result.
The solution to
˙ x = Ax + b(t) (6.7)
with initial value
x(0) = x
0
(6.8)
is
x(t) = x
h
(t) + x
p
(t) (6.9)
where
x
h
(t) = e
At
x(0) (6.10)
and
x
p
(t) =
_
t
0
e
A(t−s)
b(s) ds . (6.11)
72 CHAPTER 6. ORDINARY DIFFERENTIAL SYSTEMS
Since we now have a formula for the general solution to a linear differential system, we seem to have all we need.
However, we do not know how to compute the matrix exponential. The naive solution to use the definition (6.5)
requires too many terms for a good approximation. As we have done for the SVD and the Schur decomposition, we
will only point out that several methods exist for computing a matrix exponential, but we will not discuss how this is
done
2
. In a fundamental paper on the subject, Nineteen dubious ways to compute the exponential of a matrix (SIAM
Review, vol. 20, no. 4, pp. 801-36), Cleve Moler and Charles Van Loan discuss a large number of different methods,
pointing out that no one of them is appropriate for all situations. A full discussion of this matter is beyond the scope
of these notes.
When the matrix A is constant, as we currently assume, we can be much more specific about the structure of the
solution (6.9) of system (6.7), and particularly so about the solution x
h
(t) to the homogeneous part. Specifically, the
matrix exponential (6.10) can be written as a linear combination, with constant vector coefficients, of scalar expo-
nentials multiplied by polynomials. In the general theory of linear differential systems, this is shown via the Jordan
canonical form. However, in the paper cited above, Moler and Van Loan point out that the Jordan form cannot be
computed reliably, and small perturbations in the data can change the results dramatically. Fortunately, a similar result
can be found through the Schur decomposition introduced in chapter 5. The next section shows how to do this.
6.3 Structure of the Solution
For the homogeneous case b(t) = 0, consider the first order system of linear differential equations
˙ x = Ax (6.12)
x(0) = x
0
. (6.13)
Two cases arise: either A admits n distinct eigenvalues, or is does not. In chapter 5, we have seen that if (but not only
if) A has n distinct eigenvalues then it has n linearly independent eigenvectors (theorem 5.1.1), and we have shown
how to find x
h
(t) by solving an eigenvalue problem. In section 6.3.1, we briefly review this solution. Then, in section
6.3.2, we show how to compute the homogeneous solution x
h
(t) in the extreme case of an n × n matrix A with n
coincident eigenvalues.
To be sure, we have seen that matrices with coincident eigenvalues can still have a full set of linearly independent
eigenvectors (see for instance the identity matrix). However, the solution procedure we introduce in section 6.3.2 for
the case of n coincident eigenvalues can be applied regardless to how many linearly independent eigenvectors exist.
If the matrix has a full complement of eigenvectors, the solution obtained in section 6.3.2 is the same as would be
obtained with the method of section 6.3.1.
Once these two extreme cases (nondefective matrix or all-coincident eigenvalues) have been handled, we show a
general procedure in section 6.3.3 for solving a homogeneous or nonhomogeneous differential system for any, square,
constant matrix A, defective or not. This procedure is based on backsubstitution, and produces a result analogous
to that obtained via Jordan decomposition for the homogeneous part x
h
(t) of the solution. However, since it is
based on the numerically sound Schur decomposition, the method of section 6.3.3 is superior in practice. For a
nonhomogeneous system, the procedure can be carried out analytically if the functions in the right-hand side vector
b(t) can be integrated.
6.3.1 A is Not Defective
In chapter 5 we saw how to find the homogeneous part x
h
(t) of the solution when A has a full set of n linearly
independent eigenvectors. This result is briefly reviewed in this section for convenience.
3
If A is not defective, then it has n linearly independent eigenvectors q
1
, . . . , q
n
with corresponding eigenvalues
λ
1
, . . . , λ
n
. Let
Q =
_
q
1
· · · q
n
¸
.
2
In Matlab, expm(A) is the matrix exponential of A.
3
Parts of this subsection and of the following one are based on notes written by Scott Cohen.
6.3. STRUCTURE OF THE SOLUTION 73
This square matrix is invertible because its columns are linearly independent. Since Aq
i
= λ
i
q
i
, we have
AQ = QΛ, (6.14)
where Λ = diag(λ
1
, . . . , λ
n
) is a square diagonal matrix with the eigenvalues of A on its diagonal. Multiplying both
sides of (6.14) by Q
−1
on the right, we obtain
A = QΛQ
−1
. (6.15)
Then, system (6.12) can be rewritten as follows:
˙ x = Ax
˙ x = QΛQ
−1
x
Q
−1
˙ x = ΛQ
−1
x
˙ y = Λy, (6.16)
where y = Q
−1
x. The last equation (6.16) represents n uncoupled, homogeneous, differential equations ˙ y
i
= λ
i
y
i
.
The solution is
y
h
(t) = e
Λt
y(0),
where
e
Λt
= diag(e
λ1t
, . . . , e
λnt
).
Using the relation x = Qy, and the consequent relation y(0) = Q
−1
x(0), we see that the solution to the homogeneous
system (6.12) is
x
h
(t) = Qe
Λt
Q
−1
x(0).
If A is normal, that is, if it has n orthonormal eigenvectors q
1
, . . . q
n
, then Q is replaced by the Hermitian matrix
S =
_
s
1
· · · s
n
¸
, Q
−1
is replaced by S
H
, and the solution to (6.12) becomes
x
h
(t) = Se
Λt
S
H
x(0).
6.3.2 A Has n Coincident Eigenvalues
When A = QΛQ
−1
, we derived that the solution to (6.12) is x
h
(t) = Qe
Λt
Q
−1
x(0). Comparing with (6.6), it should
be the case that
e
Q(Λt)Q
−1
= Qe
Λt
Q
−1
.
This follows easily from the definition of e
Z
and the fact that (Q(Λt)Q
−1
)
j
= Q(Λt)
j
Q
−1
. Similarly, if A = SΛS
H
,
where S is Hermitian, then the solution to (6.12) is x
h
(t) = Se
Λt
S
H
x(0), and
e
S(Λt)S
H
= Se
Λt
S
H
.
How can we compute the matrix exponential in the extreme case in which A has n coincident eigenvalues, regard-
less of the number of its linearly independent eigenvectors? In any case, A admits a Schur decomposition
A = STS
H
(theorem 5.3.2). We recall that S is a unitary matrix and T is upper triangular with the eigenvalues of Aon its diagonal.
Thus we can write T as
T = Λ +N,
where Λ is diagonal and N is strictly upper triangular. The solution (6.6) in this case becomes
x
h
(t) = e
S(Tt)S
H
x(0) = Se
Tt
S
H
x(0) = Se
Λt+Nt
S
H
x(0).
74 CHAPTER 6. ORDINARY DIFFERENTIAL SYSTEMS
Thus we can compute (6.6) if we can compute e
Tt
= e
Λt+Nt
. This turns out to be almost as easy as computing e
Λt
when the diagonal matrix Λ is a multiple of the identity matrix:
Λ = λI
that is, when all the eigenvalues of A coincide. In fact, in this case, Λt and Nt commute:
Λt Nt = λIt Nt = λt Nt = Nt λt = Nt λIt = Nt Λt .
It can be shown that if two matrices Z
1
and Z
2
commute, that is if
Z
1
Z
2
= Z
2
Z
1
,
then
e
Z1+Z2
= e
Z1
e
Z2
= e
Z2
e
Z1
.
Thus, in our case, we can write
e
Λt+Nt
= e
Λt
e
Nt
.
We already know how to compute e
Λt
, so it remains to show how to compute e
Nt
. The fact that Nt is strictly upper
triangular makes the computation of this matrix exponential much simpler than for a general matrix Z.
Suppose, for example, that N is 4 ×4. Then N has three nonzero superdiagonals, N
2
has two nonzero superdiag-
onals, N
3
has one nonzero superdiagonal, and N
4
is the zero matrix:
N =
_
¸
¸
_
0 ∗ ∗ ∗
0 0 ∗ ∗
0 0 0 ∗
0 0 0 0
_
¸
¸
_
→N
2
=
_
¸
¸
_
0 0 ∗ ∗
0 0 0 ∗
0 0 0 0
0 0 0 0
_
¸
¸
_

N
3
=
_
¸
¸
_
0 0 0 ∗
0 0 0 0
0 0 0 0
0 0 0 0
_
¸
¸
_
→N
4
=
_
¸
¸
_
0 0 0 0
0 0 0 0
0 0 0 0
0 0 0 0
_
¸
¸
_
.
In general, for a strictly upper triangular n ×n matrix, we have N
j
= 0 for all j ≥ n (i.e., N is nilpotent of order n).
Therefore,
e
Nt
=

j=0
N
j
t
j
j!
=
n−1

j=0
N
j
t
j
j!
is simply a finite sum, and the exponential reduces to a matrix polynomial.
In summary, the general solution to the homogeneous differential system (6.12) with initial value (6.13) when the
n ×n matrix A has n coincident eigenvalues is given by
x
h
(t) = Se
Λt
n−1

j=0
N
j
t
j
j!
S
H
x
0
(6.17)
where
A = S(Λ +N)S
H
is the Schur decomposition of A,
Λ = λI
is a multiple of the identity matrix containing the coincident eigenvalues of A on its diagonal, and N is strictly upper
triangular.
6.3. STRUCTURE OF THE SOLUTION 75
6.3.3 The General Case
We are now ready to solve the linear differential system
˙ x = Ax + b(t) (6.18)
x(0) = x
0
(6.19)
in the general case of a constant matrix A, defective or not, with arbitrary b(t). In fact, let A = STS
H
be the Schur
decomposition of A, and consider the transformed system
˙ y(t) = Ty(t) + c(t) (6.20)
where
y(t) = S
H
x(t) and c(t) = S
H
b(t) . (6.21)
The triangular matrix T can always be written in the following form:
T =
_
¸
¸
¸
_
T
11
· · · · · · T
1k
0 T
22
· · · T
2k
.
.
.
.
.
.
.
.
.
0 · · · 0 T
kk
_
¸
¸
¸
_
where the diagonal blocks T
ii
for i = 1, . . . , k are of size n
i
×n
i
(possibly 1×1) and contain all-coincident eigenvalues.
The remaining nonzero blocks T
ij
with i < j can be in turn bundled into matrices
R
i
=
_
T
i,i+1
· · · T
i,k
¸
that contain everything to the right of the corresponding T
ii
. The vector c(t) can be partitioned correspondingly as
follows
c(t) =
_
¸
_
c
1
(t)
.
.
.
c
k
(t)
_
¸
_
where c
i
has n
i
entries, and the same can be done for
y(t) =
_
¸
_
y
1
(t)
.
.
.
y
k
(t)
_
¸
_
and for the initial values
y(0) =
_
¸
_
y
1
(0)
.
.
.
y
k
(0)
_
¸
_ .
The triangular system (6.20) can then be solved by backsubstitution as follows:
for i = k down to 1
if i < k
d
i
(t) = R
i
[y
i+1
(t), . . . , y
k
(t)]
T
else
d
i
(t) = 0 (an n
k
-dimensional vector of zeros)
end
T
ii
= λ
i
I +N
i
(diagonal and strictly upper-triangular part of T
ii
)
y
i
(t) = e
λ
i
It

n
i
−1
j=0
N
j
i
t
j
j!
y
i
(0) +
_
t
0
_
e
λ
i
I(t−s)

n
i
−1
j=0
N
j
i
(t−s)
j
j!
_
(c
i
(s) + d
i
(s)) ds
end.
76 CHAPTER 6. ORDINARY DIFFERENTIAL SYSTEMS
In this procedure, the expression for y
i
(t) is a direct application of equations (6.9), (6.10), (6.11), and (6.17) with
S = I. In the general case, the applicability of this routine depends on whether the integral in the expression for y
i
(t)
can be computed analytically. This is certainly the case when b(t) is a constant vector b, because then the integrand is
a linear combination of exponentials multiplied by polynomials in t −s, which can be integrated by parts.
The solution x(t) for the original system (6.18) is then
x(t) = Sy(t) .
As an illustration, we consider a very small example, the 2 ×2 homogeneous, triangular case,
_
˙ y
1
˙ y
2
_
=
_
t
11
t
12
0 t
22
_ _
y
1
y
2
_
. (6.22)
When t
11
= t
22
= λ, we obtain
y(t) = e
λIt
_
1 t
12
t
0 1
_
y(0) .
In scalar form, this becomes
y
1
(t) = (y
1
(0) +t
12
y
2
(0) t) e
λt
y
2
(t) = y
2
(0) e
λt
,
and it is easy to verify that this solution satisfies the differential system (6.22).
When t
11
= λ
1
= t
22
= λ
2
, we could solve the system by finding the eigenvectors of T, since we know that in this
case two linearly independent eigenvectors exist (theorem 5.1.1). Instead, we apply the backsubstitution procedure
introduced in this section. The second equation of the system,
˙ y
2
(t) = t
22
y
2
has solution
y
2
(t) = y
2
(0) e
λ
2
t
.
We then have
d
1
(t) = t
12
y
2
(t) = t
12
y
2
(0) e
λ2t
and
y
1
(t) = y
1
(0)e
λ1t
+
_
t
0
e
λ1(t−s)
d
1
(s) ds
= y
1
(0)e
λ1t
+t
12
y
2
(0) e
λ1t
_
t
0
e
−λ1s
e
λ2s
ds
= y
1
(0)e
λ1t
+t
12
y
2
(0) e
λ1t
_
t
0
e
(λ2−λ1)s
ds
= y
1
(0)e
λ1t
+
t
12
y
2
(0)
λ
2
−λ
1
e
λ1t
(e
(λ2−λ1)t
−1)
= y
1
(0)e
λ
1
t
+
t
12
y
2
(0)
λ
2
−λ
1
(e
λ
2
t
−e
λ
1
t
)
Exercise: verify that this solution satisfies both the differential equation (6.22) and the initial value equation y(0) = y
0
.
Thus, the solutions to system (6.22) for t
11
= t
22
and for t
11
= t
22
have different forms. While y
2
(t) is the same
in both cases, we have
y
1
(t) = y
1
(0) e
λt
+t
12
y
2
(0) t e
λt
if t
11
= t
22
y
1
(t) = y
1
(0)e
λ1t
+
t
12
y
2
(0)
λ
2
−λ
1
(e
λ2t
−e
λ1t
) if t
11
= t
22
.
6.4. A CONCRETE EXAMPLE 77
rest position
of mass 1
rest position
of mass 2
2
v
1
v
1
2
2
1
3
Figure 6.1: A system of masses and springs. In the absence of external forces, the two masses would assume the
positions indicated by the dashed lines.
This would seem to present numerical difficulties when t
11
≈ t
22
, because the solution would suddenly switch
from one form to the other as the difference between t
11
and t
22
changes from about zero to exactly zero or viceversa.
This, however, is not a problem. In fact,
lim
λ
1
→λ
e
λt
−e
λ
1
t
λ −λ
1
= t e
λt
,
and the transition between the two cases is smooth.
6.4 A Concrete Example
In this section we set up and solve a more concrete example of a system of differential equations. The initial system
has two second-order equations, and is transformed into a first-order system with four equations. The 4 ×4 matrix of
the resulting system has an interesting structure, which allows finding eigenvalues and eigenvectors analytically with a
little trick. The point of this section is to show how to transform the complex formal solution of the differential system,
computed with any of the methods described above, into a real solution in a form appropriate to the problem at hand.
Consider the mechanical system in figure 6.1. Suppose that we want to study the evolution of the system over time.
Since forces are proportional to accelerations, because of Newton’s law, and since accelerations are second derivatives
of position, the new equations are differential. Because differentiation occurs only with respect to one variable, time,
these are ordinary differential equations, as opposed to partial.
In the following we write the differential equations that describe this system. Two linear differential equations of
the second order
4
result. We will then transform these into four linear differential equations of the first order.
By Hooke’s law, the three springs exert forces that are proportional to the springs’ elongations:
f
1
= c
1
v
1
f
2
= c
2
(v
2
−v
1
)
4
Recall that the order of a differential equation is the highest degree of derivative that appears in it.
78 CHAPTER 6. ORDINARY DIFFERENTIAL SYSTEMS
f
3
= −c
3
v
2
where the c
i
are the positive spring constants (in newtons per meter).
The accelerations of masses 1 and 2 (springs are assumed to be massless) are proportional to their accelerations,
according to Newton’s second law:
m
1
¨ v
1
= −f
1
+f
2
= −c
1
v
1
+c
2
(v
2
−v
1
) = −(c
1
+c
2
)v
1
+c
2
v
2
m
2
¨ v
2
= −f
2
+f
3
= −c
2
(v
2
−v
1
) −c
3
v
2
= c
2
v
1
−(c
2
+c
3
)v
2
or, in matrix form,
¨v = Bv (6.23)
where
v =
_
v
1
v
2
_
and B =
_

c
1
+c
2
m
1
c
2
m
1
c
2
m
2

c
2
+c
3
m
2
_
.
We also assume that initial conditions
v(0) and ˙ v(0) (6.24)
are given, which specify positions and velocities of the two masses at time t = 0.
To solve the second-order system (6.23), we will first transform it to a system of four first-order equations. As
shown in the introduction to this chapter, the trick is to introduce variables to denote the first-order derivatives of v, so
that second-order derivatives of v are first-order derivatives of the new variables. For uniformity, we define four new
variables
u =
_
¸
¸
_
u
1
u
2
u
3
u
4
_
¸
¸
_
=
_
¸
¸
_
v
1
v
2
˙ v
1
˙ v
2
_
¸
¸
_
(6.25)
so that
u
3
= ˙ v
1
and u
4
= ˙ v
2
,
while the original system (6.23) becomes
_
˙ u
3
˙ u
4
_
= B
_
u
1
u
2
_
.
We can now gather these four first-order differential equations into a single system as follows:
˙ u = Au (6.26)
where
A =
_
¸
¸
_
0 0
0 0
1 0
0 1
B
0 0
0 0
_
¸
¸
_
.
Likewise, the initial conditions (6.24) are replaced by the (known) vector
u(0) =
_
v(0)
˙ v(0)
_
.
In the next section we solve equation (6.26).
6.5. SOLUTION OF THE EXAMPLE 79
6.5 Solution of the Example
Not all matrices have a full set of linearly independent eigenvectors. With the system of springs in figure 6.1, however,
we are lucky. The eigenvalues of A are solutions to the equation
Ax = λx , (6.27)
where we recall that
A =
_
0 I
B 0
_
and B =
_

c
1
+c
2
m1
c
2
m1
c
2
m2

c
2
+c
3
m2
_
.
Here, the zeros in A are 2 ×2 matrices of zeros, and I is the 2 ×2 identity matrix. If we partition the vector x into its
upper and lower halves y and z,
x =
_
y
z
_
,
we can write
Ax =
_
0 I
B 0
_ _
y
z
_
=
_
z
By
_
so that the eigenvalue equation (6.27) can be written as the following pair of equations:
z = λy (6.28)
By = λz ,
which yields
By = µy with µ = λ
2
.
In other words, the eigenvalues of A are the square roots of the eigenvalues of B: if we denote the two eigenvalues of
B as µ
1
and µ
2
, then the eigenvalues of A are
λ
1
=

µ
1
λ
2
= −

µ
1
λ
3
=

µ
2
λ
4
= −

µ
2
.
The eigenvalues µ
1
and µ
2
of B are the solutions of
det(B −µI) =
_
c
1
+c
2
m
1

__
c
2
+c
3
m
2

_

c
2
2
m
1
m
2
= µ
2
+ 2αµ +β = 0
where
α =
1
2
_
c
1
+c
2
m
1
+
c
2
+c
3
m
2
_
and β =
c
1
c
2
+c
1
c
3
+c
2
c
3
m
1
m
2
are positive constants that depend on the elastic properties of the springs and on the masses. We then obtain
µ
1,2
= −α ±γ ,
where
γ =
_
α
2
−β =
¸
1
4
_
c
1
+c
2
m
1

c
2
+c
3
m
2
_
2
+
c
2
2
m
1
m
2
.
The constant γ is real because the radicand is nonnegative. We also have that α ≥ γ, so that the two solutions µ
1,2
are
real and negative, and the four eigenvalues of A,
λ
1
=

−α +γ , λ
2
= −

−α +γ , (6.29)
λ
3
=

−α −γ , λ
4
= −

−α −γ (6.30)
80 CHAPTER 6. ORDINARY DIFFERENTIAL SYSTEMS
come in nonreal, complex-conjugate pairs. This is to be expected, since our system of springs obviously exhibits an
oscillatory behavior.
Also the eigenvectors of A can be derived from those of B. In fact, from equation (6.28) we see that if y is an
eigenvector of B corresponding to eigenvalue µ = λ
2
, then there are two corresponding eigenvectors for A of the
form
x =
_
y
±λy
_
. (6.31)
The eigenvectors of B are the solutions of
(B −(−α ±γ)I)y = 0 . (6.32)
Since ±(−α±γ) are eigenvalues of B, the determinant of this equation is zero, and the two scalar equations in (6.32)
must be linearly dependent. The first equation reads

_
c
1
+c
2
m
1
−α ±γ
_
y
1
+
c
2
m
1
y
2
= 0
and is obviously satisfied by any vector of the form
y = k
_
c
2
m
1
c
1
+c
2
m
1
−α ±γ
_
where k is an arbitrary constant. For k = 0, y denotes the two eigenvectors of B, and from equation (6.31) the four
eigenvectors of A are proportional to the four columns of the following matrix:
Q =
_
¸
¸
_
c2
m1
c2
m1
c2
m1
c2
m1
a +λ
2
1
a +λ
2
2
a +λ
2
3
a +λ
2
4
λ
1
c2
m
1
λ
2
c2
m
1
λ
3
c2
m
1
λ
4
c2
m
1
λ
1
_
a +λ
2
1
_
λ
2
_
a +λ
2
2
_
λ
3
_
a +λ
2
3
_
λ
4
_
a +λ
2
4
_
_
¸
¸
_
(6.33)
where
a =
c
1
+c
2
m
1
.
The general solution to the first-order differential system (6.26) is then given by equation (6.17). Since we just found
four distinct eigenvectors, however, we can write more simply
u(t) = Qe
Λt
Q
−1
u(0) (6.34)
where
Λ =
_
¸
¸
_
λ
1
0 0 0
0 λ
2
0 0
0 0 λ
3
0
0 0 0 λ
4
_
¸
¸
_
.
In these expressions, the values of λ
i
are given in equations (6.30), and Q is in equation (6.33).
Finally, the solution to the original, second-order system (6.23) can be obtained from equation (6.25) by noticing
that v is equal to the first two components of u.
This completes the solution of our system of differential equations. However, it may be useful to add some
algebraic manipulation in order to show that the solution is indeed oscillatory. As we see in the following, the masses’
motions can be described by the superposition of two sinusoids whose frequencies depend on the physical constants
involved (masses and spring constants). The amplitudes and phases of the sinusoids, on the other hand, depend on the
initial conditions.
To simplify our manipulation, we note that
u(t) = Qe
Λt
w ,
6.5. SOLUTION OF THE EXAMPLE 81
where we defined
w = Q
−1
u(0) . (6.35)
We now leave the constants in w unspecified, and derive the general solution v(t) for the original, second-order
problem. Numerical values for the constants can be found from the initial conditions u(0) by equation (6.35). We
have
v(t) = Q(1 : 2, :)e
Λt
w ,
where Q(1 : 2, :) denotes the first two rows of Q. Since
λ
2
= −λ
1
and λ
4
= −λ
3
(see equations (6.30)), we have
Q(1 : 2, :) =
_
q
1
q
1
q
2
q
2
¸
where we defined
q
1
=
_
c
2
m1
c
1
+c
2
m1

2
1
_
and q
2
=
_
c
2
m1
c
1
+c
2
m1

2
3
_
.
Thus, we can write
v(t) = q
1
_
k
1
e
λ
1
t
+k
2
e
−λ
1
t
_
+ q
2
_
k
3
e
λ
3
t
+k
4
e
−λ
3
t
_
.
Since the λs are imaginary but v(t) is real, the k
i
must come in complex-conjugate pairs:
k
1
= k

2
and k
3
= k

4
. (6.36)
In fact, we have
v(0) = q
1
(k
1
+k
2
) + q
2
(k
3
+k
4
)
and from the derivative
˙ v(t) = q
1
λ
1
_
k
1
e
λ1t
−k
2
e
−λ1t
_
+ q
2
λ
3
_
k
3
e
λ3t
−k
4
e
−λ3t
_
we obtain
˙ v(0) = q
1
λ
1
(k
1
−k
2
) + q
2
λ
3
(k
3
−k
4
) .
Since the vectors q
i
are independent (assuming that the mass c
2
is nonzero), this means that
k
1
+k
2
is real k
1
−k
2
is purely imaginary
k
3
+k
4
is real k
3
−k
4
is purely imaginary ,
from which equations (6.36) follow.
Finally, by using the relation
e
jx
+e
−jx
2
= cos x ,
and simple trigonometry we obtain
v(t) = q
1
A
1
cos(ω
1
t +φ
1
) + q
2
A
2
cos(ω
2
t +φ
2
)
where
ω
1
=

α −γ =
¸
¸
¸
_
1
2
(a +b) −
¸
1
4
(a −b)
2
+
c
2
2
m
1
m
2
ω
2
=

α +γ =
¸
¸
¸
_
1
2
(a +b) +
¸
1
4
(a −b)
2
+
c
2
2
m
1
m
2
82 CHAPTER 6. ORDINARY DIFFERENTIAL SYSTEMS
and
a =
c
1
+c
2
m
1
, b =
c
2
+c
3
m
2
.
Notice that these two frequencies depend only on the configuration of the system, and not on the initial conditions.
The amplitudes A
i
and phases φ
i
, on the other hand, depend on the constants k
i
as follows:
A
1
= 2|k
1
| , A
2
= 2|k
3
|
φ
1
= arctan
2
(Im(k
1
), Re(k
1
)) φ
2
= arctan
2
(Im(k
3
), Re(k
3
))
where Re, Im denote the real and imaginary part and where the two-argument function arctan
2
is defined as follows
for (x, y) = (0, 0)
arctan
2
(y, x) =
_
¸
¸
_
¸
¸
_
arctan(
y
x
) if x > 0
π + arctan(
y
x
) if x < 0
π
2
if x = 0 and y > 0

π
2
if x = 0 and y < 0
and is undefined for (x, y) = (0, 0). This function returns the arctangent of y/x (notice the order of the arguments) in
the proper quadrant, and extends the function by continuity along the y axis.
The two constants k
1
and k
3
can be found from the given initial conditions v(0) and ˙ v(0) from equations (6.35)
and (6.25).
Chapter 7
Stochastic State Estimation
Perhaps the most important part of studying a problem in robotics or vision, as well as in most other sciences, is to
determine a good model for the phenomena and events that are involved. For instance, studying manipulation requires
defining models for how a robot arm can move and for how it interacts with the world. Analyzing image motion
implies defining models for how points move in space and how this motion projects onto the image. When motion
is involved, as is very often the case, models take on frequently the form of dynamic systems. A dynamic system
is a mathematical description of a quantity that evolves over time. The theory of dynamic systems is both rich and
fascinating. Although in this chapter we will barely scratch its surface, we will consider one of its most popular and
useful aspects, the theory of state estimation, in the particular form of Kalman filtering. To this purpose, an informal
definition of a dynamic system is given in the next section. The definition is then illustrated by setting up the dynamic
system equations for a simple but realistic application, that of modeling the trajectory of an enemy mortar shell. In
sections 7.3 through 7.5, we will develop the theory of the Kalman filter, and in section 7.6 we will see that the shell
can be shot down before it hits us. As discussed in section 7.7, Kalman filtering has intimate connections with the
theory of algebraic linear systems we have developed in chapters 2 and 3.
7.1 Dynamic Systems
In its most general meaning, the term system refers to some physical entity on which some action is performed by
means of an input u. The system reacts to this input and produces an output y (see figure 7.1).
A dynamic system is a system whose phenomena occur over time. One often says that a system evolves over time.
Simple examples of a dynamic system are the following:
• An electric circuit, whose input is the current in a given branch and whose output is a voltage across a pair of
nodes.
• A chemical reactor, whose inputs are the external temperature, the temperature of the gas being supplied, and
the supply rate of the gas. The output can be the temperature of the reaction product.
• A mass suspended from a spring. The input is the force applied to the mass and the output is the position of the
mass.
input output system
S u y
Figure 7.1: A general system.
83
84 CHAPTER 7. STOCHASTIC STATE ESTIMATION
In all these examples, what is input and what is output is a choice that depends on the application. Also, all the
quantities in the examples vary continuously with time. In other cases, as for instance for switching networks and
computers, it is more natural to consider time as a discrete variable. If time varies continuously, the system is said to
be continuous; if time varies discretely, the system is said to be discrete.
7.1.1 State
Given a dynamic system, continuous or discrete, the modeling problem is to somehow correlate inputs (causes) with
outputs (effects). The examples above suggest that the output at time t cannot be determined in general by the value
assumed by the input quantity at the same point in time. Rather, the output is the result of the entire history of the
system. An effort of abstraction is therefore required, which leads to postulating a new quantity, called the state, which
summarizes information about the past and the present of the system. Specifically, the value x(t) taken by the state at
time t must be sufficient to determine the output at the same point in time. Also, knowledge of both x(t
1
) and u
[t
1
,t
2
)
,
that is, of the state at time t
1
and the input over the interval t
1
≤ t < t
2
, must allow computing the state (and hence
the output) at time t
2
. For the mass attached to a spring, for instance, the state could be the position and velocity of
the mass. In fact, the laws of classical mechanics allow computing the new position and velocity of the mass at time t
2
given its position and velocity at time t
1
and the forces applied over the interval [t
1
, t
2
). Furthermore, in this example,
the output y of the system happens to coincide with one of the two state variables, and is therefore always deducible
from the latter.
Thus, in a dynamic system the input affects the state, and the output is a function of the state. For a discrete
system, the way that the input changes the state at time instant number k into the new state at time instant k + 1 can
be represented by a simple equation:
x
k+1
= f(x
k
, u
k
, k)
where f is some function that represents the change, and u
k
is the input at time k. Similarly, the relation between state
and output can be expressed by another function:
y
k
= h(x
k
, k) .
A discrete dynamic system is completely described by these two equations and an initial state x
0
. In general, all
quantities are vectors.
For continuous systems, time does not come in quanta, so one cannot compute x
k+1
as a function of x
k
, u
k
, and
k, but rather compute x(t
2
) as a functional φ of x(t
1
) and the entire input u over the interval [t
1
, t
2
):
x(t
2
) = φ(x(t
1
), u(·), t
1
, t
2
)
where u(·) represents the entire function u, not just one of its values. A description of the system in terms of func-
tions, rather than functionals, can be given in the case of a regular system, for which the functional φ is continuous,
differentiable, and with continuous first derivative. In that case, one can show that there exists a function f such that
the state x(t) of the system satisfies the differential equation
˙ x(t) = f(x(t), u(t), t)
where the dot denotes differentiation with respect to time. The relation from state to output, on the other hand, is
essentially the same as for the discrete case:
y(t) = h(x(t), t) .
Specifying the initial state x
0
completes the definition of a continuous dynamic system.
7.1.2 Uncertainty
The systems defined in the previous section are called deterministic, since the evolution is exactly determined once
the initial state x at time 0 is known. Determinism implies that both the evolution function f and the output function
h are known exactly. This is, however, an unrealistic state of affairs. In practice, the laws that govern a given physical
7.2. AN EXAMPLE: THE MORTAR SHELL 85
system are known up to some uncertainty. In fact, the equations themselves are simple abstractions of a complex
reality. The coefficients that appear in the equations are known only approximately, and can change over time as a
result of temperature changes, component wear, and so forth. A more realistic model then allows for some inherent,
unresolvable uncertainty in both f and h. This uncertainty can be represented as noise that perturbs the equations we
have presented so far. A discrete system then takes on the following form:
x
k+1
= f(x
k
, u
k
, k) +η
k
y
k
= h(x
k
, k) +ξ
k
and for a continuous system
˙ x(t) = f(x(t), u(t), t) +η(t)
y(t) = h(x(t), t) +ξ(t) .
Without loss of generality, the noise distributions can be assumed to have zero mean, for otherwise the mean can be
incorporated into the deterministic part, that is, in either f or h. The mean may not be known, but this is a different
story: in general the parameters that enter into the definitions of f and h must be estimated by some method, and the
mean perturbations are no different.
A common assumption, which is sometimes valid and always simplifies the mathematics, is that η and ξ are
zero-mean Gaussian random variables with known covariance matrices Q and R, respectively.
7.1.3 Linearity
The mathematics becomes particularly simple when both the evolution function f and the output function h are linear.
Then, the system equations become
x
k+1
= F
k
x
k
+G
k
u
k

k
y
k
= H
k
x
k

k
for the discrete case, and
˙ x(t) = F(t)x(t) +G(t)u(t) +η(t)
y(t) = H(t)x(t) +ξ(t)
for the continuous one. It is useful to specify the sizes of the matrices involved. We assume that the input u is a vector
in R
p
, the state x is in R
n
, and the output y is in R
m
. Then, the state propagation matrix F is n ×n, the input matrix
G is n × p, and the output matrix H is m × n. The covariance matrix Q of the system noise η is n × n, and the
covariance matrix of the output noise ξ is m×m.
7.2 An Example: the Mortar Shell
In this section, the example of the mortar shell will be discussed in order to see some of the technical issues involved
in setting up the equations of a dynamic system. In particular, we consider discretization issues because the physical
system is itself continuous, but we choose to model it as a discrete system for easier implementation on a computer.
In sections 7.3 through 7.5, we consider the state estimation problem: given observations of the output y over an
interval of time, we want to determine the state x of the system. This is a very important task. For instance, in the case
of the mortar shell, the state is the (initially unknown) position and velocity of the shell, while the output is a set of
observations made by a tracking system. Estimating the state then leads to enough knowledge about the shell to allow
driving an antiaircraft gun to shoot the shell down in mid-flight.
You spotted an enemy mortar installation about thirty kilometers away, on a hill that looks about 0.5 kilometers
higher than your own position. You want to track incoming projectiles with a Kalman filter so you can aim your guns
86 CHAPTER 7. STOCHASTIC STATE ESTIMATION
accurately. You do not know the initial velocity of the projectiles, so you just guess some values: 0.6 kilometers/second
for the horizontal component, 0.1 kilometers/second for the vertical component. Thus, your estimate of the initial state
of the projectile is
ˆx
0
=
_
¸
¸
_
˙
d
d
˙ z
z
_
¸
¸
_
=
_
¸
¸
_
−0.6
30
0.1
0.5
_
¸
¸
_
where d is the horizontal coordinate, z is the vertical, you are at (0, 0), and dots denote derivatives with respect to
time.
From your high-school physics, you remember that the laws of motion for a ballistic trajectory are the following:
d(t) = d(0) +
˙
d(0)t (7.1)
z(t) = z(0) + ˙ z(0)t −
1
2
gt
2
(7.2)
where g is the gravitational acceleration, equal to 9.8 × 10
−3
kilometers per second squared. Since you do not trust
your physics much, and you have little time to get ready, you decide to ignore air drag. Because of this, you introduce
a state update covariance matrix Q = 0.1I
4
, where I
4
is the 4 ×4 identity matrix.
All you have to track the shells is a camera pointed at the mortar that will rotate so as to keep the projectile at the
center of the image, where you see a blob that increases in size as the projectile gets closer. Thus, the aiming angle of
the camera gives you elevation information about the projectile’s position, and the size of the blob tells you something
about the distance, given that you know the actual size of the projectiles used and all the camera parameters. The
projectile’s elevation is
e = 1000
z
d
(7.3)
when the projectile is at (d, z). Similarly, the size of the blob in pixels is
s =
1000

d
2
+z
2
. (7.4)
You do not have very precise estimates of the noise that corrupts e and s, so you guess measurement covariances
R
e
= R
s
= 1000, which you put along the diagonal of a 2 ×2 diagonal measurement covariance matrix R.
7.2.1 The Dynamic System Equation
Equations (7.1) and (7.2) are continuous. Since you are taking measurements every dt = 0.2 seconds, you want to
discretize these equations. For the z component, equation (7.2) yields
z(t +dt) −z(t) = z(0) + ˙ z(0)(t +dt) −
1
2
g(t +dt)
2

_
z(0) + ˙ z(0)t −
1
2
gt
2
_
= ( ˙ z(0) −gt)dt −
1
2
g(dt)
2
= ˙ z(t)dt −
1
2
g(dt)
2
,
since ˙ z(0) −gt = ˙ z(t).
Consequently, if t +dt is time instant k + 1 and t is time instant k, you have
z
k+1
= z
k
+ ˙ z
k
dt −
1
2
g(dt)
2
. (7.5)
The reasoning for the horizontal component d is the same, except that there is no acceleration:
d
k+1
= d
k
+
˙
d
k
dt . (7.6)
7.2. AN EXAMPLE: THE MORTAR SHELL 87
Equations (7.5) and (7.6) can be rewritten as a single system update equation
x
k+1
= Fx
k
+Gu
where
x
k
=
_
¸
¸
_
˙
d
k
d
k
˙ z
k
z
k
_
¸
¸
_
is the state, the 4 × 4 matrix F depends on dt, the control scalar u is equal to −g, and the 4 × 1 control matrix G
depends on dt. The two matrices F and G are as follows:
F =
_
¸
¸
_
1 0 0 0
dt 1 0 0
0 0 1 0
0 0 dt 1
_
¸
¸
_
G =
_
¸
¸
_
0
0
dt
−dt
2
2
_
¸
¸
_
.
7.2.2 The Measurement Equation
The two nonlinear equations (7.3) and (7.4) express the available measurements as a function of the true values of the
projectile coordinates d and z. We want to replace these equations with linear approximations. To this end, we develop
both equations as Taylor series around the current estimate and truncate them after the linear term. From the elevation
equation (7.3), we have
e
k
= 1000
z
k
d
k
≈ 1000
_
ˆ z
k
ˆ
d
k
+
z
k
− ˆ z
k
ˆ
d
k

ˆ z
k
ˆ
d
2
k
(d
k

ˆ
d
k
)
_
,
so that after simplifying we can redefine the measurement to be the discrepancy from the estimated value:
e

k
= e
k
−1000
ˆ z
k
ˆ
d
k
≈ 1000(
z
k
ˆ
d
k

ˆ z
k
ˆ
d
2
k
d
k
) . (7.7)
We can proceed similarly for equation (7.4):
s
k
=
1000
_
d
2
k
+z
2
k

1000
_
ˆ
d
2
k
+ ˆ z
2
k

1000
ˆ
d
k
(
ˆ
d
2
k
+ ˆ z
2
k
)
3/2
(d
k

ˆ
d
k
) −
1000ˆ z
k
(
ˆ
d
2
k
+ ˆ z
2
k
)
3/2
(z
k
− ˆ z
k
)
and after simplifying:
s

k
= s
k

2000
_
ˆ
d
2
+ ˆ z
2
≈ −1000
_
ˆ
d
k
(
ˆ
d
2
k
+ ˆ z
2
k
)
3/2
d
k
+
ˆ z
k
(
ˆ
d
2
k
+ ˆ z
2
k
)
3/2
z
k
_
. (7.8)
The two measurements s

k
and e

k
just defined can be collected into a single measurement vector
y
k
=
_
s

k
e

k
_
,
and the two approximate measurement equations (7.7) and (7.8) can be written in the matrix form
y
k
= H
k
x
k
(7.9)
where the measurement matrix H
k
depends on the current state estimate ˆx
k
:
H
k
= −1000
_
_
0
ˆ
d
k
(
ˆ
d
2
k
+ˆ z
2
k
)
3/2
0
ˆ z
k
(
ˆ
d
2
k
+ˆ z
2
k
)
3/2
0
ˆ z
k
ˆ
d
2
k
0 −
1
ˆ
d
k
_
_
88 CHAPTER 7. STOCHASTIC STATE ESTIMATION
As the shell approaches us, we frantically start studying state estimation, and in particular Kalman filtering, in the
hope to build a system that lets us shoot down the shell before it hits us. The next few sections will be read under this
impending threat.
Knowing the model for the mortar shell amounts to knowing the laws by which the object moves and those that
relate the position of the projectile to our observations. So what else is there left to do? From the observations, we
would like to know where the mortar shell is right now, and perhaps predict where it will be in a few seconds, so we
can direct an antiaircraft gun to shoot down the target. In other words, we want to know x
k
, the state of the dynamic
system. Clearly, knowing x
0
instead is equivalent, at least when the dynamics of the system are known exactly (the
system noise η
k
is zero). In fact, from x
0
we can simulate the system up until time t, thereby determining x
k
as well.
Most importantly, we do not want to have all the observations before we shoot: we would be dead by then. A scheme
that refines an initial estimation of the state as new observations are acquired is called a recursive
1
state estimation
system. The Kalman filter is one of the most versatile schemes for recursive state estimations. The original paper
by Kalman (R. E. Kalman, “A new approach to linear filtering and prediction problems,” Transactions of the ASME
Journal Basic Engineering, 82:34–45, 1960) is still one of the most readable treatments of this subject from the point
of view of stochastic estimation.
Even without noise, a single observation y
k
may not be sufficient to determine the state x
k
(in the example, one
observation happens to be sufficient). This is a very interesting aspect of state estimation. It is really the ensemble of
all observations that let one estimate the state, and yet observations are processed one at a time, as they become avail-
able. A classical example of this situation in computer vision is the reconstruction of three-dimensional shape from
a sequence of images. A single image is two-dimensional, so by itself it conveys no three-dimensional information.
Kalman filters exist that recover shape information froma sequence of images. See for instance L. Matthies, T. Kanade,
and R. Szeliski, “Kalman filter-based algorithms for estimating depth from image sequences,” International Journal of
Computer Vision, 3(3):209-236, September 1989; and T.J. Broida, S. Chandrashekhar, and R. Chellappa, “Recursive
3-D motion estimation from a monocular image sequence,” IEEE Transactions on Aerospace and Electronic Systems,
26(4):639–656, July 1990.
Here, we introduce the Kalman filter from the simpler point of view of least squares estimation, since we have
developed all the necessary tools in the first part of this course. The next section defines the state estimation problem
for a discrete dynamic system in more detail. Then, section 7.4 defines the essential notions of estimation theory
that are necessary to understand the quantitative aspects of Kalman filtering. Section 7.5 develops the equation of the
Kalman filter, and section 7.6 reconsiders the example of the mortar shell. Finally, section 7.7 establishes a connection
between the Kalman filter and the solution of a linear system.
7.3 State Estimation
In this section, the estimation problem is defined in some more detail. Given a discrete dynamic system
x
k+1
= F
k
x
k
+G
k
u
k

k
(7.10)
y
k
= H
k
x
k

k
(7.11)
where the system noise η
k
and the measurement noise ξ
k
are Gaussian variables,
η
k
∼ N(0, Q
k
)
ξ
k
∼ N(0, R
k
) ,
as well as a (possibly completely wrong) estimate ˆx
0
of the initial state and an initial covariance matrix P
0
of the
estimate ˆx
0
, the Kalman filter computes the optimal estimate ˆx
k|k
at time k given the measurements y
0
, . . . , y
k
. The
filter also computes an estimate P
k|k
of the covariance of ˆx
k|k
given those measurements. In these expressions, the hat
means that the quantity is an estimate. Also, the first k in the subscript refers to which variable is being estimated, the
second to which measurements are being used for the estimate. Thus, in general, ˆx
i|j
is the estimate of the value that
x assumes at time i given the first j + 1 measurements y
0
, . . . , y
j
.
1
The term “recursive” in the systems theory literature corresponds loosely to “incremental” or “iterative” in computer science.
7.3. STATE ESTIMATION 89
k | k-1
^
y
H
k
x
k | k-1
^
x
^
k | k
k | k
P
x
^
P
k+1 | k
k+1 | k
propagate propagate
x
^
P
k-1 | k-1
k-1 | k-1
y
k
y
k-1
y
k+1
k
update
k | k-1
P
k-1 k+1
time
Figure 7.2: The update stage of the Kalman filter changes the estimate of the current system state x
k
to make the
prediction of the measurement closer to the actual measurement y
k
. Propagation then accounts for the evolution of the
system state, as well as the consequent growing uncertainty.
7.3.1 Update
The covariance matrix P
k|k
must be computed in order to keep the Kalman filter running, in the following sense. At
time k, just before the new measurement y
k
comes in, we have an estimate ˆx
k|k−1
of the state vector x
k
based on the
previous measurements y
0
, . . . , y
k−1
. Now we face the problem of incorporating the new measurement y
k
into our
estimate, that is, of transforming ˆx
k|k−1
into ˆx
k|k
. If ˆx
k|k−1
were exact, we could compute the new measurement y
k
without even looking at it, through the measurement equation (7.11). Even if ˆx
k|k−1
is not exact, the estimate
ˆy
k|k−1
= H
k
ˆx
k|k−1
is still our best bet. Now y
k
becomes available, and we can consider the residue
r
k
= y
k
−ˆy
k|k−1
= y
k
−H
k
ˆx
k|k−1
.
If this residue is nonzero, we probably need to correct our estimate of the state x
k
, so that the new prediction
ˆy
k|k
= H
k
ˆx
k|k
of the measurement value is closer to the measurement y
k
than the old prediction
ˆy
k|k−1
= H
k
ˆx
k|k−1
that we made just before the new measurement y
k
was available.
The question however is, by how much should we correct our estimate of the state? We do not want to make ˆy
k|k
coincide with y
k
. That would mean that we trust the new measurement completely, but that we do not trust our state
estimate ˆx
k|k−1
at all, even if the latter was obtained through a large number of previous measurements. Thus, we
need some criterion for comparing the quality of the new measurement y
k
with that of our old estimate ˆx
k|k−1
of the
state. The uncertainty about the former is R
k
, the covariance of the observation error. The uncertainty about the state
just before the new measurement y
k
becomes available is P
k|k−1
. The update stage of the Kalman filter uses R
k
and
P
k|k−1
to weigh past evidence (ˆx
k|k−1
) and new observations (y
k
). This stage is represented graphically in the middle
of figure 7.2. At the same time, also the uncertainty measure P
k|k−1
must be updated, so that it becomes available for
the next step. Because a new measurement has been read, this uncertainty becomes usually smaller: P
k|k
< P
k|k−1
.
The idea is that as time goes by the uncertainty on the state decreases, while that about the measurements may
remain the same. Then, measurements count less and less as the estimate approaches its true value.
90 CHAPTER 7. STOCHASTIC STATE ESTIMATION
7.3.2 Propagation
Just after arrival of the measurement y
k
, both state estimate and state covariance matrix have been updated as described
above. But between time k and time k + 1 both state and covariance may change. The state changes according to the
system equation (7.10), so our estimate ˆx
k+1|k
of x
k+1
given y
0
, . . . , y
k
should reflect this change as well. Similarly,
because of the system noise η
k
, our uncertainty about this estimate may be somewhat greater than one time epoch ago.
The system equation (7.10) essentially “dead reckons” the new state from the old, and inaccuracies in our model of
how this happens lead to greater uncertainty. This increase in uncertainty depends on the system noise covariance Q
k
.
Thus, both state estimate and covariance must be propagated to the new time k + 1 to yield the new state estimate
ˆx
k+1|k
and the new covariance P
k+1|k
. Both these changes are shown on the right in figure 7.2.
In summary, just as the state vector x
k
represents all the information necessary to describe the evolution of a
deterministic system, the covariance matrix P
k|k
contains all the necessary information about the probabilistic part of
the system, that is, about how both the system noise η
k
and the measurement noise ξ
k
corrupt the quality of the state
estimate ˆx
k|k
.
Hopefully, this intuitive introduction to Kalman filtering gives you an idea of what the filter does, and what infor-
mation it needs to keep working. To turn these concepts into a quantitative algorithm we need some preliminaries on
optimal estimation, which are discussed in the next section. The Kalman filter itself is derived in section 7.5.
7.4 BLUE Estimators
In what sense does the Kalman filter use covariance information to produce better estimates of the state? As we will
se later, the Kalman filter computes the Best Linear Unbiased Estimate (BLUE) of the state. In this section, we see
what this means, starting with the definition of a linear estimation problem, and then considering the attributes “best”
and “unbiased” in turn.
7.4.1 Linear Estimation
Given a quantity y (the observation) that is a known function of another (deterministic but unknown) quantity x (the
state) plus some amount of noise,
y = h(x) + n , (7.12)
the estimation problem amounts to finding a function
ˆx = L(y)
such that ˆx is as close as possible to x. The function L is called an estimator, and its value ˆx given the observations y
is called an estimate. Inverting a function is an example of estimation. If the function h is invertible and the noise term
n is zero, then L is the inverse of h, no matter how the phrase “as close as possible” is interpreted. In fact, in that case
ˆx is equal to x, and any distance between ˆx and x must be zero. In particular, solving a square, nonsingular system
y = Hx (7.13)
is, in this somewhat trivial sense, a problem of estimation. The optimal estimator is then represented by the matrix
L = H
−1
and the optimal estimate is
ˆx = Ly .
A less trivial example occurs, for a linear observation function, when the matrix H has more rows than columns,
so that the system (7.13) is overconstrained. In this case, there is usually no inverse to H, and again one must say in
what sense ˆx is required to be “as close as possible” to x. For linear systems, we have so far considered the criterion
that prefers a particular ˆx if it makes the Euclidean norm of the vector y − Hx as small as possible. This is the
7.4. BLUE ESTIMATORS 91
(unweighted) least squares criterion. In section 7.4.2, we will see that in a very precise sense ordinary least squares
solve a particular type of estimation problem, namely, the estimation problem for the observation equation (7.12) with
h a linear function and n Gaussian zero-mean noise with the indentity matrix for covariance.
An estimator is said to be linear if the function L is linear. Notice that the observation function h can still be
nonlinear. If L is required to be linear but h is not, we will probably have an estimator that produces a worse estimate
than a nonlinear one. However, it still makes sense to look for the best possible linear estimator. The best estimator
for a linear observation function happens to be a linear estimator.
7.4.2 Best
In order to define what is meant by a “best” estimator, one needs to define a measure of goodness of an estimate. In
the least squares approach to solving a linear system like (7.13), this distance is defined as the Euclidean norm of the
residue vector
y −Hˆx
between the left and the right-hand sides of equation (7.13), evaluated at the solution ˆx. Replacing (7.13) by a “noisy
equation”,
y = Hx + n (7.14)
does not change the nature of the problem. Even equation (7.13) has no exact solution when there are more independent
equations than unknowns, so requiring equality is hopeless. What the least squares approach is really saying is that
even at the solution ˆx there is some residue
n = y −Hˆx (7.15)
and we would like to make that residue as small as possible in the sense of the Euclidean norm. Thus, an overcon-
strained system of the form (7.13) and its “noisy” version (7.14) are really the same problem. In fact, (7.14) is the
correct version, if the equality sign is to be taken literally.
The noise term, however, can be used to generalize the problem. In fact, the Euclidean norm of the residue (7.15)
treats all components (all equations in (7.14)) equally. In other words, each equation counts the same when computing
the norm of the residue. However, different equations can have noise terms of different variance. This amounts to
saying that we have reasons to prefer the quality of some equations over others or, alternatively, that we want to
enforce different equations to different degrees. From the point of view of least squares, this can be enforced by some
scaling of the entries of n or, even, by some linear transformation of them:
n →Wn
so instead of minimizing n
2
= n
T
n (the square is of course irrelevant when it comes to minimization), we now
minimize
Wn
2
= n
T
R
−1
n
where
R
−1
= W
T
W
is a symmetric, nonnegative-definite matrix. This minimization problem, called weighted least squares, is only slightly
different from its unweighted version. In fact, we have
Wn = W(y −Hx) = Wy −WHx
so we are simply solving the system
Wy = WHx
in the traditional, “unweighted” sense. We know the solution from normal equations:
ˆx = ((WH)
T
WH)
−1
(WH)
T
Wy = (H
T
R
−1
H)
−1
H
T
R
−1
y .
92 CHAPTER 7. STOCHASTIC STATE ESTIMATION
Interestingly, this same solution is obtained from a completely different criterion of goodness of a solution ˆx. This
criterion is a probabilistic one. We consider this different approach because it will let us show that the Kalman filter is
optimal in a very useful sense.
The new criterion is the so-called minimum-covariance criterion. The estimate ˆx of x is some function of the
measurements y, which in turn are corrupted by noise. Thus, ˆx is a function of a random vector (noise), and is therefore
a random vector itself. Intuitively, if we estimate the same quantity many times, from measurements corrupted by
different noise samples from the same distribution, we obtain different estimates. In this sense, the estimates are
random.
It makes therefore sense to measure the quality of an estimator by requiring that its variance be as small as possible:
the fluctuations of the estimate ˆx with respect to the true (unknown) value x from one estimation experiment to the
next should be as small as possible. Formally, we want to choose a linear estimator L such that the estimates ˆx = Ly
it produces minimize the following covariance matrix:
P = E[(x −ˆx)(x −ˆx)
T
] .
Minimizing a matrix, however, requires a notion of “size” for matrices: how large is P? Fortunately, most inter-
esting matrix norms are equivalent, in the sense that given two different definitions P
1
and P
2
of matrix norm
there exist two positive scalars α, β such that
αP
1
< P
2
< βP
1
.
Thus, we can pick any norm we like. In fact, in the derivations that follow, we only use properties shared by all norms,
so which norm we actually use is irrelevant. Some matrix norms were mentioned in section 3.2.
7.4.3 Unbiased
In additionto requiring our estimator to be linear and with minimum covariance, we also want it to be unbiased, in the
sense that if repeat the same estimation experiment many times we neither consistently overestimate nor consistently
underestimate x. Mathematically, this translates into the following requirement:
E[x −ˆx] = 0 and E[ˆx] = E[x] .
7.4.4 The BLUE
We now address the problem of finding the Best Linear Unbiased Estimator (BLUE)
ˆx = Ly
of x given that y depends on x according to the model (7.14), which is repeated here for convenience:
y = Hx + n . (7.16)
First, we give a necessary and sufficient condition for L to be unbiased.
Lemma 7.4.1 Let n in equation (7.16) be zero mean. Then the linear estimator L is unbiased if an only if
LH = I ,
the identity matrix.
Proof.
E[x −ˆx] = E[x −Ly] = E[x −L(Hx + n)]
= E[(I −LH)x] −E[Ln] = (I −HL)E[x]
7.4. BLUE ESTIMATORS 93
since E[Ln] = LE[n] and E[n] = 0. For this to hold for all x we need I −LH = 0. ∆
And now the main result.
Theorem 7.4.2 The Best Linear Unbiased Estimator (BLUE)
ˆx = Ly
for the measurement model
y = Hx + n
where the noise vector n has zero mean and covariance R is given by
L = (H
T
R
−1
H)
−1
H
T
R
−1
and the covariance of the estimate ˆx is
P = E[(x −ˆx)(x −ˆx)
T
] = (H
T
R
−1
H)
−1
. (7.17)
Proof. We can write
P = E[(x −ˆx)(x −ˆx)
T
] = E[(x −Ly)(x −Ly)
T
]
= E[(x −LHx −Ln)(x −LHx −Ln)
T
] = E[((I −LH)x −Ln)((I −LH)x −Ln)
T
]
= E[Lnn
T
L
T
] = LE[nn
T
] L
T
= LRL
T
because L is unbiased, so that LH = I.
To show that
L
0
= (H
T
R
−1
H)
−1
H
T
R
−1
(7.18)
is the best choice, let L be any (other) linear unbiased estimator. We can trivially write
L = L
0
+ (L −L
0
)
and
P = LRL
T
= [L
0
+ (L −L
0
)]R[L
0
+ (L −L
0
)]
T
= L
0
RL
T
0
+ (L −L
0
)RL
T
0
+L
0
R(L −L
0
)
T
+ (L −L
0
)R(L −L
0
)
T
.
From (7.18) we obtain
RL
T
0
= RR
−1
H(H
T
R
−1
H)
−1
= H(H
T
R
−1
H)
−1
so that
(L −L
0
)RL
T
0
= (L −L
0
)H(H
T
R
−1
H)
−1
= (LH −L
0
H)(H
T
R
−1
H)
−1
.
But L and L
0
are unbiased, so LH = L
0
H = I, and
(L −L
0
)RL
T
0
= 0 .
The term L
0
R(L −L
0
)
T
is the transpose of this, so it is zero as well. In conclusion,
P = L
0
RL
T
0
+ (L −L
0
)R(L −L
0
)
T
,
the sum of two positive definite or at least semidefinite matrices. For such matrices, the norm of the sum is greater or
equal to either norm, so this expression is minimized when the second term vanishes, that is, when L = L
0
.
This proves that the estimator given by (7.18) is the best, that is, that it has minimum covariance. To prove that the
covariance P of ˆx is given by equation (7.17), we simply substitute L
0
for L in P = LRL
T
:
P = L
0
RL
T
0
= (H
T
R
−1
H)
−1
H
T
R
−1
RR
−1
H(H
T
R
−1
H)
−1
= (H
T
R
−1
H)
−1
H
T
R
−1
H(H
T
R
−1
H)
−1
= (H
T
R
−1
H)
−1
as promised. ∆
94 CHAPTER 7. STOCHASTIC STATE ESTIMATION
7.5 The Kalman Filter: Derivation
We now have all the components necessary to write the equations for the Kalman filter. To summarize, given a linear
measurement equation
y = Hx + n
where n is a Gaussian random vector with zero mean and covariance matrix R,
n ∼ N(0, R) ,
the best linear unbiased estimate ˆx of x is
ˆx = PH
T
R
−1
y
where the matrix
P

= E[(ˆx −x)(ˆx −x)
T
] = (H
T
R
−1
H)
−1
is the covariance of the estimation error.
Given a dynamic system with system and measurement equations
x
k+1
= F
k
x
k
+G
k
u
k

k
(7.19)
y
k
= H
k
x
k

k
where the system noise η
k
and the measurement noise ξ
k
are Gaussian random vectors,
η
k
∼ N(0, Q
k
)
ξ
k
∼ N(0, R
k
) ,
as well as the best, linear, unbiased estimate ˆx
0
of the initial state with an error covariance matrix P
0
, the Kalman
filter computes the best, linear, unbiased estimate ˆx
k|k
at time k given the measurements y
0
, . . . , y
k
. The filter also
computes the covariance P
k|k
of the error ˆx
k|k
−x
k
given those measurements. Computation occurs according to the
phases of update and propagation illustrated in figure 7.2. We now apply the results from optimal estimation to the
problem of updating and propagating the state estimates and their error covariances.
7.5.1 Update
At time k, two pieces of data are available. One is the estimate ˆx
k|k−1
of the state x
k
given measurements up to but not
including y
k
. This estimate comes with its covariance matrix P
k|k−1
. Another way of saying this is that the estimate
ˆx
k|k−1
differs from the true state x
k
by an error term e
k
whose covariance is P
k|k−1
:
ˆx
k|k−1
= x
k
+ e
k
(7.20)
with
E[e
k
e
T
k
] = P
k|k−1
.
The other piece of data is the new measurement y
k
itself, which is related to the state x
k
by the equation
y
k
= H
k
x
k

k
(7.21)
with error covariance
E[ξ
k
ξ
T
k
] = R
k
.
We can summarize this available information by grouping equations 7.20 and 7.21 into one, and packaging the error
covariances into a single, block-diagonal matrix. Thus, we have
y = Hx
k
+ n
7.5. THE KALMAN FILTER: DERIVATION 95
where
y =
_
ˆx
k|k−1
y
k
_
, H =
_
I
H
k
_
, n =
_
e
k
ξ
k
_
,
and where n has covariance
R =
_
P
k|k−1
0
0 R
k
_
.
As we know, the solution to this classical estimation problem is
ˆx
k|k
= P
k|k
H
T
R
−1
y
P
k|k
= (H
T
R
−1
H)
−1
.
This pair of equations represents the update stage of the Kalman filter. These expressions are somewhat wasteful,
because the matrices H and R contain many zeros. For this reason, these two update equations are now rewritten in a
more efficient and more familiar form. We have
P
−1
k|k
= H
T
R
−1
H
=
_
I H
T
k
¸
_
P
−1
k|k−1
0
0 R
−1
k
_ _
I
H
k
_
= P
−1
k|k−1
+H
T
k
R
−1
k
H
k
and
ˆx
k|k
= P
k|k
H
T
R
−1
y
= P
k|k
_
P
−1
k|k−1
H
T
k
R
−1
k
_
_
ˆx
k|k−1
y
k
_
= P
k|k
(P
−1
k|k−1
ˆx
k|k−1
+H
T
k
R
−1
k
y
k
)
= P
k|k
((P
−1
k|k
−H
T
k
R
−1
k
H
k
)ˆx
k|k−1
+H
T
k
R
−1
k
y
k
)
= ˆx
k|k−1
+P
k|k
H
T
k
R
−1
k
(y
k
−H
k
ˆx
k|k−1
) .
In the last line, the difference
r
k

= y
k
−H
k
ˆx
k|k−1
is the residue between the actual measurement y
k
and its best estimate based on ˆx
k|k−1
, and the matrix
K
k

= P
k|k
H
T
k
R
−1
k
is usually referred to as the Kalman gain matrix, because it specifies the amount by which the residue must be multi-
plied (or amplified) to obtain the correction term that transforms the old estimate ˆx
k|k−1
of the state x
k
into its new
estimate ˆx
k|k
.
7.5.2 Propagation
Propagation is even simpler. Since the new state is related to the old through the system equation 7.19, and the noise
term η
k
is zero mean, unbiasedness requires
ˆx
k+1|k
= F
k
ˆx
k|k
+G
k
u
k
,
96 CHAPTER 7. STOCHASTIC STATE ESTIMATION
which is the state estimate propagation equation of the Kalman filter. The error covariance matrix is easily propagated
thanks to the linearity of the expectation operator:
P
k+1|k
= E[(ˆx
k+1|k
−x
k+1
)(ˆx
k+1|k
−x
k+1
)
T
]
= E[(F
k
(ˆx
k|k
−x
k
) −η
k
)(F
k
(ˆx
k|k
−x
k
) −η
k
)
T
]
= F
k
E[(ˆx
k|k
−x
k
)(ˆx
k|k
−x
k
)
T
]F
T
k
+E[η
k
η
T
k
]
= F
k
P
k|k
F
T
k
+Q
k
where the system noise η
k
and the previous estimation error ˆx
k|k
−x
k
were assumed to be uncorrelated.
7.5.3 Kalman Filter Equations
In summary, the Kalman filter evolves an initial estimate and an initial error covariance matrix,
ˆx
0|−1

= ˆx
0
and P
0|−1

= P
0
,
both assumed to be given, by the update equations
ˆx
k|k
= ˆx
k|k−1
+K
k
(y
k
−H
k
ˆx
k|k−1
)
P
−1
k|k
= P
−1
k|k−1
+H
T
k
R
−1
k
H
k
where the Kalman gain is defined as
K
k
= P
k|k
H
T
k
R
−1
k
and by the propagation equations
ˆx
k+1|k
= F
k
ˆx
k|k
+G
k
u
k
P
k+1|k
= F
k
P
k|k
F
T
k
+Q
k
.
7.6 Results of the Mortar Shell Experiment
In section 7.2, the dynamic system equations for a mortar shell were set up. Matlab routines available through the
class Web page implement a Kalman filter (with naive numerics) to estimate the state of that system from simulated
observations. Figure 7.3 shows the true and estimated trajectories. Notice that coincidence of the trajectories does not
imply that the state estimate is up-to-date. For this it is also necessary that any given point of the trajectory is reached
by the estimate at the same time instant. Figure 7.4 shows that the distance between estimated and true target position
does indeed converge to zero, and this occurs in time for the shell to be shot down. Figure 7.5 shows the 2-norm of the
covariance matrix over time. Notice that the covariance goes to zero only asymptotically.
7.7 Linear Systems and the Kalman Filter
In order to connect the theory of state estimation with what we have learned so far about linear systems, we now show
that estimating the initial state x
0
from the first k +1 measurements, that is, obtaining ˆx
0|k
, amounts to solving a linear
system of equations with suitable weights for its rows.
The basic recurrence equations (7.10) and (7.11) can be expanded as follows:
y
k
= H
k
x
k

k
= H
k
(F
k−1
x
k−1
+G
k−1
u
k−1

k−1
) +ξ
k
= H
k
F
k−1
x
k−1
+H
k
G
k−1
u
k−1
+H
k
η
k−1

k
= H
k
F
k−1
(F
k−2
x
k−2
+G
k−2
u
k−2

k−2
) +H
k
G
k−1
u
k−1
+H
k
η
k−1

k
7.7. LINEAR SYSTEMS AND THE KALMAN FILTER 97
5 10 15 20 25 30 35
−0.2
0
0.2
0.4
0.6
0.8
1
1.2
1.4
true (dashed) and estimated (solid) missile trajectory
Figure 7.3: The true and estimated trajectories get closer to one another. Trajectories start on the right.
0 5 10 15 20 25 30
0
0.2
0.4
0.6
0.8
1
1.2
1.4
1.6
1.8
2
distance between true and estimated missile position vs. time
Figure 7.4: The estimate actually closes in towards the target.
98 CHAPTER 7. STOCHASTIC STATE ESTIMATION
0 5 10 15 20 25 30
0
5
10
15
20
25
30
35
40
norm of the state covariance matrix vs time
Figure 7.5: After an initial increase in uncertainty, the norm of the state covariance matrix converges to zero. Upwards
segments correspond to state propagation, downwards ones to state update.
= H
k
F
k−1
F
k−2
x
k−2
+H
k
(F
k−1
G
k−2
u
k−2
+G
k−1
u
k−1
) +
H
k
(F
k−1
η
k−2

k−1
) +ξ
k
.
.
.
= H
k
F
k−1
. . . F
0
x
0
+H
k
(F
k−1
. . . F
1
G
0
u
0
+. . . +G
k−1
u
k−1
) +
H
k
(F
k−1
. . . F
1
η
0
+. . . +η
k−1
) +ξ
k
or in a more compact form,
y
k
= H
k
Φ(k −1, 0)x
0
+H
k
k

j=1
Φ(k −1, j)G
j−1
u
j−1

k
(7.22)
where
Φ(l, j) =
_
F
l
. . . F
j
for l ≥ j
1 for l < j
and the term
ν
k
= H
k
k

j=1
Φ(k −1, j)η
j−1

k
is noise.
The key thing to notice about this somewhat intimidating expression is that for any k it is a linear system in x
0
, the
initial state of the system. We can write one system like the one in equation (7.22) for every value of k = 0, . . . , K,
where K is the last time instant considered, and we obtain a large system of the form
z
K
= Ψ
K
x
0
+ g
K
+ n
K
(7.23)
where
z
K
=
_
¸
_
y
0
.
.
.
y
K
_
¸
_
7.7. LINEAR SYSTEMS AND THE KALMAN FILTER 99
Ψ
K
=
_
¸
¸
¸
_
H
0
H
1
F
0
.
.
.
H
K
Φ(K −1, 0)
_
¸
¸
¸
_
g
K
=
_
¸
¸
¸
_
0
H
1
G
0
u
0
.
.
.
H
K
(Φ(K −1, 1)G
0
u
0
+. . . + Φ(K −1, K)G
K−1
u
K−1
)
_
¸
¸
¸
_
n
K
=
_
¸
_
ν
0
.
.
.
ν
K
_
¸
_ .
Without knowing anything about the statistics of the noise vector n
K
in equation (7.23), the best we can do is to
solve the system
z
K
= Ψ
K
x
0
+ g
K
in the sense of least squares, to obtain an estimate of x
0
from the measurements y
0
, . . . , y
K
:
ˆx
0|K
= Ψ

K
(z
K
−g
K
)
where Ψ

K
is the pseudoinverse of Ψ
K
. We know that if Ψ
K
has full rank, the result with the pseudoinverse is the
same as we would obtain by solving the normal equations, so that
Ψ

K
= (Ψ
T
K
Ψ
K
)
−1
Ψ
T
K
.
The least square solution to system (7.23) minimizes the residue between the left and the right-hand side under the
assumption that all equations are to be treated the same way. This is equivalent to assuming that all the noise terms in
n
K
are equally important. However, we know the covariance matrices of all these noise terms, so we ought to be able
to do better, and weigh each equation to keep these covariances into account. Intuitively, a small covariance means that
we believe in that measurement, and therefore in that equation, which should consequently be weighed more heavily
than others. The quantitative embodiment of this intuitive idea is at the core of the Kalman filter.
In summary, the Kalman filter for a linear system has been shown to be equivalent to a linear equation solver, under
the assumption that the noise that affects each of the equations has the same probability distribution, that is, that all
the noise terms in n
K
in equation 7.23 are equally important. However, the Kalman filter differs from a linear solver
in the following important respects:
1. The noise terms in n
K
in equation 7.23 are not equally important. Measurements come with covariance matrices,
and the Kalman filter makes optimal use of this information for a proper weighting of each of the scalar equations
in (7.23). Better information ought to yield more accurate results, and this is in fact the case.
2. The system (7.23) is not solved all at once. Rather, an initial solution is refined over time as new measurements
become available. The final solution can be proven to be exactly equal to solving system (7.23) all at once.
However, having better and better approximations to the solution as new data come in is much preferable in a
dynamic setting, where one cannot in general wait for all the data to be collected. In some applications, data my
never stop arriving.
3. A solution for the estimate ˆx
k|k
of the current state is given, and not only for the estimate ˆx
0|k
of the initial state.
As time goes by, knowledge of the initial state may obsolesce and become less and less useful. The Kalman
filter computes up-to-date information about the current state.

2

Chapter 1

Introduction
Fields such as robotics or computer vision are interdisciplinary subjects at the intersection of engineering and computer science. By their nature, they deal with both computers and the physical world. Although the former are in the latter, the workings of computers are best described in the black-and-white vocabulary of discrete mathematics, which is foreign to most classical models of reality, quantum physics notwithstanding. This class surveys some of the key tools of applied math to be used at the interface of continuous and discrete. It is not on robotics or computer vision, nor does it cover any other application area. Applications evolve rapidly, but their mathematical foundations remain. Even if you will not pursue any of these fields, the mathematics that you learn in this class will not go wasted. To be sure, applied mathematics is a discipline in itself and, in many universities, a separate department. Consequently, this class can be a quick tour at best. It does not replace calculus or linear algebra, which are assumed as prerequisites, nor is it a comprehensive survey of applied mathematics. What is covered is a compromise between the time available and what is useful and fun to talk about. Even if in some cases you may have to wait until you take an applied class to fully appreciate the usefulness of a particular topic, I hope that you will enjoy studying these subjects in their own right.

1.1 Who Should Take This Class
The main goal of this class is to present a collection of mathematical tools for both understanding and solving problems in fields that manipulate models of the real world, such as robotics, artificial intelligence, vision, engineering, or several aspects of the biological sciences. Several classes at most universities each cover some of the topics presented in this class, and do so in much greater detail. If you want to understand the full details of any one of the topics in the syllabus below, you should take one or more of these other classes instead. If you want to understand how these tools are implemented numerically, you should take one of the classes in the scientific computing program, which again cover these issues in much better detail. Finally, if you want to understand robotics, vision, or other applied fields, you should take classes in these subjects, since this course is not on applications. On the other hand, if you do plan to study robotics, vision, or other applied subjects in the future, and you regard yourself as a user of the mathematical techniques outlined in the syllabus below, then you may benefit from this course. Of the proofs, we will only see those that add understanding. Of the implementation aspects of algorithms that are available in, say, Matlab or LApack, we will only see the parts that we need to understand when we use the code. In brief, we will be able to cover more topics than other classes because we will be often (but not always) unconcerned with rigorous proof or implementation issues. The emphasis will be on intuition and on practicality of the various algorithms. For instance, why are singular values important, and how do they relate to eigenvalues? What are the dangers of Newton-style minimization? How does a Kalman filter work, and why do PDEs lead to sparse linear systems? In this spirit, for instance, we discuss Singular Value Decomposition and Schur decomposition both because they never fail and because they clarify the structure of an algebraic or a differential linear problem.

3

1.3 2.2 Partial differential equations and sparse linear systems 4.1.1.5 4.3 4.2.1.3 4.2 4.3 The Kalman filter 4.3 Constraints and Lagrange multipliers 3.2.4 3.2 Levenberg-Marquardt method 2.2.2 Function optimization 2.3.1 Linear estimation 3.2.2 3. symmetric.1.2 4. div.2.2 Weighted least squares 3.1. Unknown functions of several variables 4.1.2 Syllabus Here is the ideal syllabus.2 Statistical estimation 3. surface.4 CHAPTER 1. INTRODUCTION 1.1 4. Introduction 2. Unknown functions of one real variable 3.2.1. Unknown numbers 2.2 2.3 Calculus of variations 4.4 4.1. curl Line.2.1 Newton and Gauss-Newton methods 2. and volume integrals Green’s theorem and potential fields of two variables Stokes’ and divergence theorems and potential fields of three variables Diffusion and flow problems Finite differences Direct versus iterative solution methods Jacobi and Gauss-Seidel iterations Successive overrelaxation 4.1 Algebraic linear systems 2.4 Grad. positive-definite matrices 3.3 3.2.1 Ordinary differential linear systems 3.1 3.1.1.1 Euler-Lagrange equations 4. 1.3.1.5 Eigenvalues and eigenvectors The Schur decomposition Ordinary differential linear systems The matrix zoo Real.1 2.2.2.4 Characterization of the solutions to a linear system Gaussian elimination The Singular Value Decomposition The pseudoinverse 2. but how much we cover depends on how fast we go.2 The brachistochrone .1.1 Tensor fields of several variables 4.1 4.1.

1.3. DISCUSSION OF THE SYLLABUS

5

1.3 Discussion of the Syllabus
In robotics, vision, physics, and any other branch of science whose subject belongs to or interacts with the real world, mathematical models are developed that describe the relationship between different quantities. Some of these quantities are measured, or sensed, while others are inferred by calculation. For instance, in computer vision, equations tie the coordinates of points in space to the coordinates of corresponding points in different images. Image points are data, world points are unknowns to be computed. Similarly, in robotics, a robot arm is modeled by equations that describe where each link of the robot is as a function of the configuration of the link’s own joints and that of the links that support it. The desired position of the end effector, as well as the current configuration of all the joints, are the data. The unknowns are the motions to be imparted to the joints so that the end effector reaches the desired target position. Of course, what is data and what is unknown depends on the problem. For instance, the vision system mentioned above could be looking at the robot arm. Then, the robot’s end effector position could be the unknowns to be solved for by the vision system. Once vision has solved its problem, it could feed the robot’s end-effector position as data for the robot controller to use in its own motion planning problem. Sensed data are invariably noisy, because sensors have inherent limitations of accuracy, precision, resolution, and repeatability. Consequently, the systems of equations to be solved are typically overconstrained: there are more equations than unknowns, and it is hoped that the errors that affect the coefficients of one equation are partially cancelled by opposite errors in other equations. This is the basis of optimization problems: Rather than solving a minimal system exactly, an optimization problem tries to solve many equations simultaneously, each of them only approximately, but collectively as well as possible, according to some global criterion. Least squares is perhaps the most popular such criterion, and we will devote a good deal of attention to it. In summary, the problems encountered in robotics and vision, as well as other applications of mathematics, are optimization problems. A fundamental distinction between different classes of problems reflects the complexity of the unknowns. In the simplest case, unknowns are scalars. When there is more than one scalar, the unknown is a vector of numbers, typically either real or complex. Accordingly, the first part of this course will be devoted to describing systems of algebraic equations, especially linear equations, and optimization techniques for problems whose solution is a vector of reals. The main tool for understanding linear algebraic systems is the Singular Value Decomposition (SVD), which is both conceptually fundamental and practically of extreme usefulness. When the systems are nonlinear, they can be solved by various techniques of function optimization, of which we will consider the basic aspects. Since physical quantities often evolve over time, many problems arise in which the unknowns are themselves functions of time. This is our second class of problems. Again, problems can be cast as a set of equations to be solved exactly, and this leads to the theory of Ordinary Differential Equations (ODEs). Here, “ordinary” expresses the fact that the unknown functions depend on just one variable (e.g., time). The main conceptual tool for addressing ODEs is the theory of eigenvalues, and the primary computational tool is the Schur decomposition. Alternatively, problems with time varying solutions can be stated as minimization problems. When viewed globally, these minimization problems lead to the calculus of variations. When the minimization problems above are studied locally, they become state estimation problems, and the relevant theory is that of dynamic systems and Kalman filtering. The third category of problems concerns unknown functions of more than one variable. The images taken by a moving camera, for instance, are functions of time and space, and so are the unknown quantities that one can compute from the images, such as the distance of points in the world from the camera. This leads to Partial Differential equations (PDEs), or to extensions of the calculus of variations. In this class, we will see how PDEs arise, and how they can be solved numerically.

6

CHAPTER 1. INTRODUCTION

1.4 Books
The class will be based on these lecture notes, and additional notes handed out when necessary. Other useful references include the following. • R. Courant and D. Hilbert, Methods of Mathematical Physics, Volume I and II, John Wiley and Sons, 1989. • D. A. Danielson, Vectors and Tensors in Engineering and Physics, Addison-Wesley, 1992. • J. W. Demmel, Applied Numerical Linear Algebra, SIAM, 1997. • A. Gelb et al., Applied Optimal Estimation, MIT Press, 1974. • P. E. Gill, W. Murray, and M. H. Wright, Practical Optimization, Academic Press, 1993. • G. H. Golub and C. F. Van Loan, Matrix Computations, 2nd Edition, Johns Hopkins University Press, 1989, or 3rd edition, 1997. • W. H. Press, B. P. Flannery, S. A. Teukolsky, and W. T. Vetterling, Numerical Recipes in C, 2nd Edition, Cambridge University Press, 1992. • G. Strang, Introduction to Applied Mathematics, Wellesley- Cambridge Press, 1986. • A. E. Taylor and W. R. Mann, Advanced Calculus, 3rd Edition, John Wiley and Sons, 1983. • L. N. Trefethen and D. Bau, III, Numerical Linear Algebra, SIAM, 1997. • R. Weinstock, Calculus of Variations, Dover, 1974.

Chapter 2

Algebraic Linear Systems
An algebraic linear system is a set of m equations in n unknown scalars, which appear linearly. Without loss of generality, an algebraic linear system can be written as follows: Ax = b (2.1)

where A is an m × n matrix, x is an n-dimensional vector that collects all of the unknowns, and b is a known vector of dimension m. In this chapter, we only consider the cases in which the entries of A, b, and x are real numbers. Two reasons are usually offered for the importance of linear systems. The first is apparently deep, and refers to the principle of superposition of effects. For instance, in dynamics, superposition of forces states that if force f1 produces acceleration a1 (both possibly vectors) and force f2 produces acceleration a2 , then the combined force f1 + αf2 produces acceleration a1 + αa2 . This is Newton’s second law of dynamics, although in a formulation less common than the equivalent f = ma. Because Newton’s laws are at the basis of the entire edifice of Mechanics, linearity appears to be a fundamental principle of Nature. However, like all physical laws, Newton’s second law is an abstraction, and ignores viscosity, friction, turbulence, and other nonlinear effects. Linearity, then, is perhaps more in the physicist’s mind than in reality: if nonlinear effects can be ignored, physical phenomena are linear! A more pragmatic explanation is that linear systems are the only ones we know how to solve in general. This argument, which is apparently more shallow than the previous one, is actually rather important. Here is why. Given two algebraic equations in two variables, f (x, y) = g(x, y) = we can eliminate, say, y and obtain the equivalent system F (x) = 0 y = h(x) . Thus, the original system is as hard to solve as it is to find the roots of the polynomial F in a single variable. Unfortunately, if f and g have degrees df and dg , the polynomial F has generically degree df dg . Thus, the degree of a system of equations is, roughly speaking, the product of the degrees. For instance, a system of m quadratic equations corresponds to a polynomial of degree 2m . The only case in which the exponential is harmless is when its base is 1, that is, when the system is linear. In this chapter, we first review a few basic facts about vectors in sections 2.1 through 2.4. More specifically, we develop enough language to talk about linear systems and their solutions in geometric terms. In contrast with the promise made in the introduction, these sections contain quite a few proofs. This is because a large part of the course material is based on these notions, so we want to make sure that the foundations are sound. In addition, some of the proofs lead to useful algorithms, and some others prove rather surprising facts. Then, in section 2.5, we characterize the solutions of linear algebraic systems. 7 0 0,

bm  (2. xn . . Equation (2. . .8 CHAPTER 2. . . j=k xj aj xk (2.  . Make sure that you are comfortable with this. .” . j=1 (2.2) is said to be a linear combination of a1 . the columns of a matrix A are dependent if there is a nonzero solution to the homogeneous system (2. .5). not all of which are zero. j=k 1 “iff” x j aj means “if and only if.1 The vectors a1 . .3) For later reference. This is important. We have n n xj aj = xk ak + j=1 j=1.  . . xn .  x1  . dependency means that there is a nonzero vector x such that n xj aj = 0 . . .2) is the same as Ax = b and equation (2.  b= . . Vectors that are not dependent are independent. .3) is the same as Ax = 0 where A= a1 ··· an . . . take the time to write out the components of each expression for a small example. . . such that n xj aj = 0 . . . . . .  x= . . j=k n xj aj = 0 so that ak = − j=1. Proof. In one direction.1. The vectors a1 .4) If you are not convinced of these equivalences. they are linearly dependent if there is a set of coefficients x1 . . Theorem 2. . Thus.6) as desired. the vector n b= j=1 x j aj (2. ALGEBRAIC LINEAR SYSTEMS 2. . . The converse is proven similarly: if n ak = j=1. an with coefficients x1 . . In other words. xn   b1  . an and n real numbers x1 . j=1 Let xk be nonzero for some k. an are linearly dependent iff1 at least one of them is a linear combination of the others. an are linearly dependent if they admit the null vector as a nonzero linear combination. it is useful to rewrite the last two equalities in a different form.5) (2. . .1 Linear (In)dependence Given n vectors a1 . xn .

Proof. . this implies that every linear space contains the zero vector. 0=b−b= n n n xj aj − j=1 j=1 xj aj = j=1 (xj − xj )aj ∆ but because the aj are linearly independent. . and state the following Lemma 2. . Theorem 2.6).1. .2 Basis A set a1 . . . BASIS for some k. . the coefficients x1 . . . an of A are linearly independent and the system Ax = b admits a solution. then the solution is unique. . Let also b= j=1 n xj aj . Just let k be the last of the nonzero xj . xn such that n b= j=1 xj aj are uniquely determined.2. . . .1 Given a vector b in the vector space B and a basis a1 . which then becomes n ak = j<k xj aj xk ∆ as desired. ∆2 We can make the first part of the proof above even more specific. an for B.2. An equivalent formulation is the following: If the columns a1 .2 If n nonzero vectors a1 . . B is said to be a vector space if it contains all the linear combinations of its basis vectors. . . . The previous theorem is a very important result. then 9 n xj aj = 0 j=1 by letting xk = −1 (so that x is nonzero). an are linearly dependent then at least one of them is a linear combination of the ones that precede it. The basis vectors are said to span the vector space. . this is possible only when xj − xj = 0 for every j. Proof. Then xj = 0 for j > k in (2. . Then. 2 This symbol marks the end of a proof. .2. an is said to be a basis for a set B of vectors if the aj are linearly independent and every vector in B can be written as a linear combination of them. 2. . In particular.

Since all bases for a space have the same number of vectors. . CHAPTER 2. So it must be one of the aj vectors. This vector cannot be a1 .2 implies that any other basis for Rm has m vectors. This shows that n ≥ n. We call a set of vectors that contains a basis for B a generating set for B. That is. since any vector  b1  . since the removed vector depends on the others. In fact. . .2 Two different bases for the same vector space B have the same number of vectors. . G is still a generating set for B. an be two different bases for B. Since B is a vector space.2. since the vectors in it cannot generate the left-over a s. the vectors of the set G = a1 . writing it right after a1 : G = a1 .  b= . . it makes sense to define the dimension of a space as the number of vectors in any of its bases. and that there are still left-over a vectors that have not been put into G. . From lemma 2. Removing the latter keeps G a generating set. ALGEBRAIC LINEAR SYSTEMS Theorem 2. because at every step we have made sure that G remains a generating set. they are mutually linearly independent.2. .  . This proves that n = n . Then we show that we cannot run out of a vectors before we run out of a vectors. Since this elementary basis has m vectors. the basis of elementary vectors ej = jth column of the m × m identity matrix is clearly a basis for Rm . which proves that n ≥ n . . bm m  can be written as b= bj e j j=1 and the ej are clearly independent. ∆ A consequence of this theorem is that any basis for Rm has m vectors. We then switch the roles of a and a vectors to conclude that n ≥ n. all the a s are in B. . . .1. . The a vectors cannot run out first. Let a1 . . a2 . an . and can therefore be written as a linear combination of a1 . . G is a generating set for B. a1 . Proof. . Then each aj is in B (why?). . Since the a s form a basis. Let us continue this procedure until we run out of either a vectors to remove or a vectors to add. But then G cannot be a generating set. we must run out of a s first (or simultaneously with the last a). theorem 2. This is absurd. Another consequence of theorem 2. Suppose in fact per absurdum that G is now made only of a vectors. .2. . since it has no other vectors preceding it. Now we can add a2 to G. . . Consequently. an and a1 . and the two results together imply that n = n . which are independent of those in G. an must be linearly dependent. since any basis for Rm can only have m vectors. Now we can repeat the whole procedure with the roles of a vectors and a vectors exchanged. Consequently. The rest of the proof now proceeds as follows: we keep removing a vectors from G and replacing them with a vectors in such a way as to keep G a generating set for B.10 Pause for a minute to verify that this formulation is equivalent. Thus. one of the vectors in G is a linear combination of those preceding it. n ≥ n .2. .2 is that n vectors of dimension m < n are bound to be dependent.

8) Thus.1). In conclusion.2.3 Inner Product and Orthogonality In this section we establish the geometric meaning of the algebraic notions of norm.1: The law of cosines states that a2 = b2 + c2 − 2bc cos θ. In the previous section we saw that any vector in Rm can be written as the linear combination m b= j=1 bj e j (2. j This result extends Pythagoras’ theorem to m dimensions. Also. and orthogonality. the square of the length b of b is m b 2 = b2 + 1 j=2 bj e j 2 . then b 2 = bT b . any two of these vectors form a 90-degree angle. How long is b? From equation (2. By Pythagoras’ theorem.7) we obtain m b = b1 e1 + j=2 bj ej and the two vectors b1 e1 and m j=2 bj ej are orthogonal. (2. . and the symbol T makes them into row vectors.7) of the elementary vectors that point along the coordinate axes. A special case of this law is Pythagoras’ theorem. b. inner product.3. Pythagoras’ theorem can now be applied again to the last sum by singling out its first term b2 e2 . projection. and so forth. vectors are column vectors by default. c (see figure 2. The length of these elementary vectors is clearly one. INNER PRODUCT AND ORTHOGONALITY 11 2. The fundamental geometric fact that is assumed to be known is the law of cosines: given a triangle with sides a. the squared length of a vector is the inner product of the vector with itself. because each of them goes from the origin to the unit point of one of the axes. obtained when θ = ±π/2. m b 2 = j=1 b2 . If we define the inner product of two m-dimensional vectors as follows: m b c= j=1 T bj cj . c a θ b Figure 2. we have a2 = b2 + c2 − 2bc cos θ where θ is the angle between the sides of length b and c. Here and elsewhere. because the coordinate axes are orthogonal by construction.

the projection vector p has the form p = ac.3. where a is some real number.3.2.8) we obtain bT b + cT c − 2bT c = bT b + cT c − 2 b Canceling equal terms and dividing by -2 yields the desired result. See figure 2.3. The line from the endpoint of b to p is orthogonal to c. the previous theorem yields bT c = 0. cT c Proof. the line between p and the endpoint of b is shortest when it is orthogonal to c: cT (b − ac) = 0 . When θ = ±π/2. ∆ Corollary 2. CHAPTER 2.2: The vector from the origin to point p is the projection of b onto c. c cos θ . the projection of b onto c is the vector from the origin to the point p on the line through c that is nearest to the endpoint of b.3 The projection of b onto c is the vector p = Pc b where Pc is the following square matrix: Pc = ccT . b p c Figure 2. ALGEBRAIC LINEAR SYSTEMS c cos θ The law of cosines applied to the triangle with sides b . From elementary geometry.2 Two nonzero vectors b and c in Rm are mutually orthogonal iff bT c = 0.12 Theorem 2. Proof. Since by definition point p is on the line through c. c . Theorem 2. Proof. and b − c yields b−c 2 = b 2 + c 2 −2 b c cos θ and from equation (2. ∆ Given two vectors b and c applied to the origin.1 bT c = b where θ is the angle between b and c.

. It is important to understand exactly what is being mapped into what in order to determine whether a linear system has solutions. . ∆ Theorem 2. . an . and its rank.4. In the process.2. If vector space A is a subspace of Rm for some m. . . Proof. an for a subspace A of Rm can be extended into a basis for Rm by adding m − n vectors an+1 .2 (Gram-Schmidt) Given n vectors a1 . With these tools. . and if so how many. we will be able to characterize the solutions to a linear system in section 2. . we also introduce a useful procedure (Gram-Schmidt) for orthonormalizing a set of linearly independent vectors. then the orthogonal complement of A is the set of all vectors in Rm that are orthogonal to all the vectors in A. This section introduces the notion of orthogonality between spaces. Two vector spaces A and B are said to be orthogonal to one another when every vector in A is orthogonal to every vector in B.1 Any basis a1 . . the complement of the xy plane in R3 is all of R3 except the xy plane. an . the following construction r=0 for j = 1 to n aj = aj − if aj = 0 r =r+1 a qr = aj j end end r T l=1 (ql aj )ql yields a set of orthonormal 3 vectors q1 . If n = m we are done.4. the given basis cannot generate all of Rm . . . .4. ORTHOGONAL SUBSPACES AND THE RANK OF A MATRIX which yields a= so that p = ac = c a = as advertised. call it an+1 . This argument can be repeated until the basis spans all of Rm . For instance. . qr that span the same space as a1 . Proof. until m = n. If r = 1. Theorem 2. . there is little to prove. . am . Let us now assume that the procedure 3 Orthonormal means orthogonal and with unit norm.4 Orthogonal Subspaces and the Rank of a Matrix Linear transformations map spaces into spaces. . . an . .5. . . . while the orthogonal complement of the xy plane is the z axis. Notice that complement and orthogonal complement are very different notions. cT b cT c ccT b cT c 13 ∆ 2. so there must be a vector. . We first prove by induction on r that the vectors qr are mutually orthonormal. The normalization in the above procedure ensures that q1 has unit norm. . If n < m. that is linearly independent of a1 . that is. defines the null space and range of a matrix. . .

This is the point of the following theorem. . . r − 1. ALGEBRAIC LINEAR SYSTEMS above has been performed a number j − 1 of times sufficient to find r − 1 vectors q1 . In symbols. then dim(A) + dim(A⊥ ) = m . . and equal in number to the number of linearly independent vectors in a1 . Proof.1). q1 . qn span A. must be some linear combination of v1 . if computed. qT qr = 0 for i i i = 1. qm .4 The number r of linearly independent columns of any m × n matrix A is equal to the number of its independent rows. Then we can show that the n − h vectors Avh+1 . that is. ∆ We can now start to talk about matrices in terms of the subspaces associated with them. . our basis for Rn : n x= j=1 cj vj . then the space generated by the rows of A has dimension r = n − h. vn for Rn . . The null space null(A) of an m × n matrix A is the space of all n-dimensional vectors that are orthogonal to the rows of A. . . . .4. and because qT aj = 0.14 CHAPTER 2. there is an n-tuple x such that Ax = b. A has n − h linearly independent rows. it follwos that qr is orthogonal to all its predecessors. vh be a basis for null(A). . . an be a basis for A. . and r =n−h where h = dim(null(A)). .1) into a basis v1 . From theorem 2. . It is not obvious that the space generated by the columns of A has also dimension r = n − h. . . Let v1 . and that these vectors are orthonormal (the inductive assumption). qr−1 .4. these n − h vectors generate the range of A. . Proof. In fact. . . . the dimension of this orthogonal space cannot exceed m − n. . .4. . . . .3 If A is a subspace of Rm and A⊥ is the orthogonal complement of A in Rm . . the dimension of the orthogonal space A⊥ is exactly m − n. .4. Avn are a basis for the range of A. ∆ qT aj i (qT aj )qT qi i i Theorem 2. vn . . . . . am for Rm (theorem 2. Because of the explicit normalization step qr = aj / aj .3. . an . . because otherwise we would have more than m vectors in a basis for Rm . . by constructing a basis for range(A). On the other hand. . .4. x ∈ null(A) iff Ax = 0.4. Now we show that the number of independent columns is also n − h. and the inner products i T qi ql are zero by the inductive assumption. . given an arbitrary vector b ∈ range(A). Orthonormalize this basis by the Gram-Schmidt procedure (theorem 2.2) to obtain q1 . . Finally. there must be a linear combination of the columns of A that is equal to b. First. . if null(A) has dimension h. . . By construction. the vector qr . Thus. and extend this basis (theorem 2. . are orthonormal (and therefore independent). Extend this basis to a basis a1 . . has unit norm. . all vectors generated by qn+1 . qn . . . and b ∈ range(A) iff Ax = b for some x. The range of A is the space of all m-dimensional vectors that are generated by the columns of A. we notice that the vectors qj span the same space as the aj s. . because the former are linear combinations of the latter. Theorem 2. Thus. so there is a space of dimension at least m − n that is orthogonal to A. . . Then for any i < r we have r−1 qT aj = qT aj − i i l=1 (qT aj )qT ql = 0 l i because the term cancels the i-th term of the sum (remember that qT qi = 1). The n-tuple x itself. Let a1 . . being an element of Rn . . . qm are orthogonal to all vectors generated by q1 . as promised. Because the new basis is orthonormal. We have already proven that the number of independent rows is n − h.

there must exist coefficients x1 . This proves that the n − h vectors Avh+1 . we can define the rank of A to be equivalently the number of linearly independent columns or of linearly independent rows of A: rank(A) = dim(range(A)) = n − dim(null(A)) . . .1 The matrix A transforms a vector x in its null space into the zero vector. such that n xj Avj = 0 j=h+1 so that A n xj vj = 0 . . . vn are linearly independent. . dimension m − r null(A)⊥ . . . vh are a basis for null(A). . .5. Thanks to this theorem. Then there exist numbers xh+1 . ∆ in conflict with the assumption that the vectors v1 . . Avn generate range(A). dimension h range(A)⊥ . . 2. and an arbitrary vector x into a vector in range(A). h.2.5 The Solutions of a Linear System Thanks to the results of the previous sections. . b = Ax = A j=1 15 n n n cj vj = j=1 cj Avj = j=h+1 cj Avj since v1 . vh span null(A). xh such that n h n xj vj = j=h+1 j=1 xj vj . . . Theorem 2. Since the vectors v1 . . dimension r = n − h . It should be obvious from the meaning of these spaces that null(A)⊥ range(A)⊥ = range(AT ) = null(AT ) where AT is the transpose of A. that they are not. . . . . and null(A)⊥ is called the rowspace of A. . defined as the matrix obtained by exchanging the rows of A with its columns. . . dimension r = rank(A) null(A). . . we prove that the n − h vectors Avh+1 . Avn are linearly independent. . not all zero. THE SOLUTIONS OF A LINEAR SYSTEM Thus. we now have a complete picture of the four spaces associated with an m × n matrix A of rank r and null-space dimension h: range(A). . xn . Second.5. . The space range(A)⊥ is called the left nullspace of the matrix. . j=h+1 But then the vector j=h+1 xj vj is in the null space of A. . A frequently used synonym for “range” is column space. . per absurdum. . . . so that Avj = 0 for j = 1. Suppose.

but is called “underdetermined” because there are fewer (r) independent equations than there are unknowns (see next item). Other solution techniques exist for linear systems. there are solutions x to Ax = b.g. while direct methods. More on this in chapter 3. but since b is in the range of A there is a linear combination of the columns (a vector x) that produces b. since the columns of A span all of Rn .9) 4 Notice that the technical meaning of “redundant” has a stronger meaning than “with more equations than unknowns. it also allows determining the rank of a matrix. Which method to use depends on the size and structure (e. Gaussian elimination. Ay = 0 ⇒ A(x + y) = b and this generates the ∞h = ∞n−r solutions mentioned above. n: • When r = n = m. In the first case above. Consider the m × n system Ax = b (2.3. admits a solution if b ∈ range(A). has more equations (m) than unknowns (n). m.16 CHAPTER 2. the system is invertible. Let Ax = b be an m × n system (m can be less than. no matter whether the system is invertible or not. are done in a finite amount of time that can be bounded given only the size of a matrix. iterative methods solve systems in a time that depends on the accuracy required. but then for any y ∈ null(A) also x + y is a solution: Ax = b . b. 2. • When r = n and m > n. In other words. so the first two cases exhaust the possibilities for r = n. “redundant” means “with exactly one solution and with more equations than unknowns. not on b.. Most notably. there can be no linear combination of the columns (no x vector) that gives b. like Gaussian elimination. and particularly its version called reduction to echelon form is such a method. In the second.” . the equations are compatible. Also. let r = rank(A) be the number of linearly independent rows or columns of A.” The case r < n < m is possible. equal to. Then. Since b is assumed to be in the range of A. listing all possibilities does not provide an operational method for determining the type of linear system for a given pair A. Also. the system is redundant. b ∈ range(A) b ∈ range(A) ⇒ no solutions ⇒ ∞n−r solutions with the convention that ∞0 = 1. Here. compatible case. depending on the relative sizes of r. In addition to always yielding a solution. sparsity) of the matrix. and is summarized in the next section. three possibilities occur. and the system is said to be incompatible. and there is a space of dimension h = n − r of vectors x such that Ax = 0. There are more equations than unknowns. Notice that if r = n then n cannot possibly exceed m. Of course. r cannot exceed either m or n.. This means that there is exactly one x that satisfies the system. it has dimension h > 0). or greater than n). ∞k is the cardinality of a k-dimensional vector space. Notice that invertibility depends only on A. whether more information is required about the matrix of the system. All the cases are summarized in figure 2. This means that the null space is nontrivial (i.e. ALGEBRAIC LINEAR SYSTEMS This allows characterizing the set of solutions to linear system as follows. 4 • When r < n the system is underdetermined. Thus. and on numerical considerations. and exactly one solution exists.6 Gaussian Elimination Gaussian elimination is an important technique for solving linear systems.

and right-hand side. Let l be one such value of i6 . The k-th step is applied to rows k. In short.3: Types of linear systems. m. exchange rows l and k of U (k) and of c(k) . Triangularization The new entry ukp is nonzero. j of U (k) : Skip no-pivot columns If uip is zero for every i = k. by solving the transformed system Ux = c . invertible. . . so equality is preserved and solving the final system yields the same solution as solving the original one. or underdetermined. . For i = k + 1.1 Reduction to Echelon Form The matrix A is reduced to echelon form by a process in m − 1 steps. . if the matrix is square and if all columns have a pivot. . This means that • Nonzero rows precede rows with all zeros. there are no restrictions on the system. If p exceeds n stop.6. of a row. m. and preserves the equality of left. Gaussian elimination replaces the rows of this system by linear combinations of the rows themselves until A is changed into a matrix U that is in the so-called echelon form. This zeros all the entries in the column below the pivot. Here is step k. which is transformed to a new vector c. If l = k. 17 no incompatible no underdetermined which can be square or rectangular. redundant. . the “pivot column index” p is set to one. 2. When this process is finished. GAUSSIAN ELIMINATION b in range(A) yes r=n yes m=n yes invertible no redundant Figure 2. • Below each pivot is a column of zeros. we compute the solution x by backsubstitution. Initially. if any. where uij denotes entry i. . and is called the pivot. that is. Selecting the largest entry in the column leads to better round-off properties. and subtract entry k of c(k) multiplied by uip /ukp from entry i of c(k) . The last step produces U (m) = U and c(m) = c. Once the system is transformed into echelon form. .6. U is in echelon form. The first nonzero entry. The first step is applied to U (1) = A and c(1) = b. ways of selecting l here lead to different numerical properties of the algorithm. subtract row k of U (k) multiplied by uip /ukp from row i of U (k) . . means that the entire algorithm is finished.5 Row exchange Now p ≤ n and uip is nonzero for some k ≤ i ≤ m. .2. then U is upper-triangular. . 6 Different 5 “Stop” . is called a pivot. . In particular. • Each pivot lies to the right of the pivot in the row above. then increment p by 1. incompatible. m of U (k) and c(k) and produces U (k+1) and c(k+1) . The same operations are applied to the rows of A and to those of b.

Here is an equivalent solution procedure that makes this unnecessary. . If r < m. First. . = cm = 0. . It is also the rank of A.6. if any. the vector vi for i = 1.. . . solutions exist. To see this. the general solution can be obtained by running backsubstitution n − r + 1 times. If on the other hand r = m (obviously r cannot exceed m). however.11) where the xji are the free variables.   . the last m − r equations yield a subsystem of the following form:     cr+1 0  .2 Backsubstitution A system Ux = c (2. = . and n − r times for the homogeneous system. Similarly. + xjn−r vn−r (2. jr . and they can be determined by backsubstitution. and can therefore be obtained by replacing each free variable by zero during backsubstitution. . . by solving the equations starting from the last one and replacing the result in the equations higher up. or if there is one it is compatible.  . . Backsubstitutions works as follows. remove the residual system. with suitable values of the free variables. and can therefore be written in the form x = v0 + xj1 v1 + .11).10) satisfied. . If there is a residual system (i. Since no vector x can satisfy these equations. . .10) in echelon form is easily solved for x. . it is inconvenient to carry around nonnumeric symbol names (the free variables). Say that the pivot columns are j1 . . 7 An n l=ji +1  uil xl  affine function is a linear function plus a constant. that is. . leaving undetermined variables specified by their name. The vector v0 is the solution when all free variables are zero. because A and U admit exactly the same solutions and are equal in size. the presence of free variables leads to an infinity of solutions. . that is.e. ALGEBRAIC LINEAR SYSTEMS 2. r < m) and some of cr+1 . The final result is a solution with as many free parameters as there are free variables. once for the nonhomogeneous system. because they are of the form 0 = ci with ci = 0. The solution obtained by backsubstitution is an affine function7 of the free variables. 0 cm Let us call this the residual subsystem. and then transform this solution procedure into one that can be more readily implemented numerically. Then symbolic backsubstitution consists of the following sequence: for i = r downto 1  1  xji = ci − uiji end This is called symbolic backsubstitution because no numerical values are assigned to free variables. cm are nonzero.   . we first solve the system symbolically. This yields the solution in the form (2. When solving a system in echelon form numerically. Then. . cr+1 = . Let us now assume that either there is no residual system. n − r can be obtained by solving the homogeneous system Ux = 0 with xji = 1 and all other free variables equal to zero. . call the variables corresponding to the r columns with pivots the basic variables.18 CHAPTER 2. . In this system. there is no residual subsystem. . and therefore of A. then the equations corresponding to these nonzero entries are incompatible. free variables are specified by name rather than by value. the linear system admits no solutions: it is incompatible. . Notice that the vectors v1 . We are left with an r × n system. Since this is the number of independent rows of U . vn−r form a basis for the null space of U . . r is the rank of U . Let r be the index of the last nonzero row of U . . Whenever they appear in the expressions for the basic variables. . Since any value given to the free variables leaves the equality of system (2. and call the other n − r the free variables. In conclusion.

−1  3 1 so the general solution is   −2 − 3x2 − x4 −2    0 x2 = x= 1    1 1 − 3 x4 0 x4 8 Selecting the largest entry in the column at or below row k is a frequent choice. 3. row 2 multiplied by 6/3 is subtracted from row 3 for both U (2) and c(2) to yield     1 3 3 2 1 U = U (3) =  0 0 3 1  . 5  1 =A= 2 −1 3 6 −3 c(1) Reduction to echelon form transforms A and b as follows. c = c(3) =  3  . no row exchange is necessary. The residual system is 0 = 0 (compatible). corresponding to the columns with pivots.6. In 8 other words. we choose a trivial pivot selection rule: we pick the (1) first nonzero entry at or below row k in the pivot column. x2 and x4 . and this would have caused rows 1 and 2 to be switched. Throughout this example. This yields the reduced system 1 3 0 0 3 3 2 1 x= 1 3 (2. GAUSSIAN ELIMINATION 19 2. for both p = 1 and p = 2. so p is set to 3: the second pivot column is column 3. the residual subsystem is first deleted. In symbolic backsubstitution. c(2) =  3  . so no row exchange is necessary.3 An Example An example will clarify both the reduction to echelon form and backsubstitution. When applied to both U (1) and c(1) this yields     1 3 3 2 1 U (2) =  0 0 3 1  . this means that u11 = a11 = 1 is the pivot. Consider the system Ax = b where U (1)   3 2 9 5  . In the first step (k = 1).2. so the system is underdetermined. The triangularization step subtracts row 1 multiplied by 2/1 from row 2. For k = 1. are free.12) (2) (2) The basic variables are x1 and x3 .6. with ∞n−r = ∞2 solutions. 3 0  1 =b= 5  . the number of nonzero rows. there are no no-pivot columns. The other two variables. and u23 is nonzero. In the triangularization step. and r < n = 4. and the rank of U and that of A is r = 2. so the pivot column index p stays at 1. . Backsubstitution applied first to row 2 and then to row 1 yields the following expressions for the pivot variables: x3 x1 1 1 1 (c2 − u24 x4 ) = (3 − x4 ) = 1 − x4 u23 3 3 1 1 = (c1 − u12 x2 − u13 x3 − u14 x4 ) = (1 − 3x2 − 3x3 − 2x4 ) u11 1 = 1 − 3x2 − (3 − x4 ) − 2x4 = −2 − 3x2 − x4 =     −3      + x2  1  + x4    0   0   −1 0   . 0 0 0 0 0 There is one zero row in the left-hand side. 0 0 6 2 6 Notice that now (k = 2) the entries uip are zero for i = 2. and subtracts row 1 multiplied by -1/1 from row 3.

In the next chapter. Gaussian elimination is a direct method. solving the nonzero part of U x = 0 with x2 = 0 and x4 = 1 leads to x3 x1 so that = = 1 1 (−1 · 1) = − 3 3 1 1 (−3 · 0 − 3 · − 1 3  −1  0  v2 =  1   −  3 1   −2  0     x = v0 + x2 v1 + x4 v2 =   1  + x2  0    −3  1   + x4   0  0  −1 0   −1  3 1  − 2 · 1) = −1 and just as before. in the sense that the answer can be found in a number of steps that depends only on the size of the matrix A. As mentioned at the beginning of this section.20 CHAPTER 2. we study a different .12) with x2 = x4 = 0 by numerical backsubstitution yields x3 x1 so that = = 1 (3 − 1 · 0) = 1 3 1 (1 − 3 · 0 − 3 · 1 − 2 · 0) = −2 1  −2  0   v0 =   1  . 0  Finally. ALGEBRAIC LINEAR SYSTEMS This same solution can be found by the numerical backsubstitution method as follows. 0  Then v1 is found by solving the nonzero part (first two rows) of U x = 0 with x2 = 1 and x4 = 0 to obtain x3 x1 so that = = 1 (−1 · 0) = 0 3 1 (−3 · 1 − 3 · 0 − 2 · 0) = −3 1  −3  1   v1 =   0  . Solving the reduced system (2.

Finally. Singular values also allow defining a notion of approximate rank which is very useful in a large number of applications. This is an iterative method. For instance. which provide great insight into the behavior of the linear transformation represented by the matrix A. This state of affairs would seem to favor Gaussian elimination over the SVD. the number of floating point operations necessary to compute the SVD of an m × n matrix is amn2 + bn3 where a. on a regular workstation. called the singular values. very efficient algorithms for the SVD exist. More generally. b are small numbers that depend on the details of the algorithm. It also allows finding approximate solutions when the linear system in question is incompatible. the computation of the SVD is numerically well behaved. much more so than Gaussian elimination. the latter yields a much more complete answer.2. based on the so-called the Singular Value Decomposition (SVD). GAUSSIAN ELIMINATION 21 method. since it computes bases for all the four spaces mentioned above. However. In addition. meaning that an exact solution usually requires an infinite number of steps. one can compute several thousand SVDs of 5 × 5 matrices in one second. as well as a set of quantities. and the number of steps necessary to find an approximate solution depends on the desired number of correct digits. for reasons that will become apparent in the next chapter.6. .

ALGEBRAIC LINEAR SYSTEMS .22 CHAPTER 2.

. . we use these new concepts to introduce the all-important concept of the Singular Value Decomposition (SVD). Then for any i = 1. qn such that p = q1 v1 + . we make this statement more specific by showing how unit vectors1 in the rowspace are transformed by matrices.2. vn be an orthonormal basis for S. .Chapter 3 The Singular Value Decomposition In section 2. and vectors in its null space into the zero vector. . This is important. . If the coordinates of P in Rm are collected in an m-dimensional vector   p1  . pm and since P is in S. In other words. in section 3. . + qn vn = V q where V = v1 ··· vn is an m × n matrix that collects the basis for S as its columns. . we saw that a matrix transforms vectors in its domain into vectors in its range (column space). we first need to introduce the notion of orthogonal matrices. . We do this in section 3. This describes the action that a matrix has on the magnitudes of vectors as well. Then. . into the orthogonal complement of the range. there must exist coefficients   q1  . Consider a point P in S. and let v1 . 3. To this end. and may need emphasis: 1 Vectors with unit norm. .  p= .  q= .  . No nonzero vector is mapped into the left null space.1 Orthogonal Matrices Let S be an n-dimensional subspace of Rm (so that we necessarily have n ≤ m). In this section. n we have n n vT p = vT i i j=1 qj vj = j=1 q j vT vj = q i . 23 .1. i since the vj are orthonormal. that is. .  . The chapter concludes with some basic applications and examples. it must be possible to write p as a linear combination of the vj s. and interpret them geometrically as transformations between systems of orthonormal coordinates.

However. because the m × m identity matrix has m linearly independent 2 columns. j In matrix form. a matrix is orthogonal if its columns are orthonormal. and is equal to the transpose of V . (3. .1. this result is still valid when V is m × m and has orthonormal columns. so that the solution L to equation (3. .1 The left inverse of an orthogonal m × n matrix V with m ≥ n exists and is equal to the transpose of V: V TV = I .3) still holds.4) is unique. THE SINGULAR VALUE DECOMPOSITION n p= j=1 qj vj and the vectors of the basis v1 . V is certainly full rank. the distinction between left and right inverse vanishes. In fact. because it is made of orthonormal columns.24 If CHAPTER 3. since equation (3. the matrix V −1 = V T is also the right inverse of V : V square 2 Nay. . Since the matrix V V T contains the inner products between the rows of V (just as V T V is formed by the inner products of its columns). However. Then L = L(V R) = (LV )R = R. while the columns of V R are linear combinations of the n columns of V .1) q = V Tp .2) V TV = I where I is the n × n identity matrix. ⇒ V −1 V = V T V = V V −1 = V V T = I . A matrix V that satisfies equation (3. . orthonormal. so V R can have at most n linearly independent columns.3) is said to be orthogonal. for square. We can summarize this discussion as follows: Theorem 3. Thus. vn are orthonormal. Also. full-rank matrices (r = m = n). so the left and the right inverse are the same. we can collect the n equations vT vj = i into the following matrix equation: 1 0 if i = j otherwise (3. then the coefficients qj are the signed magnitudes of the projections of p onto the basis vectors: qj = vT p . V T is both the left and the right inverse: V TV = V V T = I . (3. for square orthogonal matrices.3) 2 (3. Of course. In particular. and V T is then simply said to be the inverse of V : V T = V −1 .4) comparison with equation (3. if m = n. Since the left inverse of a matrix V is defined as the matrix L such that LV = I . the argument above shows that the rows of a square orthogonal matrix are orthonormal as well.3) shows that the left inverse of an orthogonal matrix V exists. Thus. Notice that V R = I cannot possibly have a solution when m > n. Of course. suppose that there exist matrices L and R such that LV = I and V R = I. . this argument requires V to be full rank.

1. or if you stand still and the whole universe spins counterclockwise around you. and the vector b − p is orthogonal to c. The derivative of this squared distance with respect to x is the vector 2x − 2U T b 3 At least geometrically. to the columns of V ). Also.3 Let U be an orthogonal matrix. as the following theorem shows.2) causes confusion. Theorem 3. when m = n.3. Because of orthogonality. In the interpretation given above. Proof. so b−p 2 = bT b + xT x − 2bT U x .3).1. ORTHOGONAL MATRICES 25 Sometimes.3 we defined the projection p of a vector b onto another vector c as the point on the line through c that is closest to b. the difference vector between b and its projection p onto range(U ) is orthogonal to range(U ): U T (b − p) = 0 . Proof. is relativity.3 Consistently with either of these geometric interpretations. and the underlying reference frame is changed from the elementary vectors ej (that is.3. Then the matrix U U T projects any vector b onto range(U ). from the columns of I) to the vectors vj (that is. A point p in range(U ) is a linear combination of the columns of U : p = Ux where x is the vector of coefficients (as many coefficients as there are columns in U ).1. and should not be surprising: If you spin clockwise on your feet. Vx 2 = xT V T V x = xT x = x 2 . In section 2. The squared distance between b and p is b − p 2 = (b − p)T (b − p) = bT b + pT p − 2bT p = bT b + xT U T U x − 2bT U x . of point P with coordinates p into a different point Q with coordinates q. however. Alternatively. This. One solution may be more efficient than the other in other ways. U T U is the identity matrix. ∆ We conclude this section with an obvious but useful consequence of orthogonality. we have the following result: Theorem 3. Furthermore.2 The norm of a vector x is not changed by multiplication by an orthogonal matrix V : Vx = x . the projection matrix is ccT (theorem 2. An analogous result holds for subspace projection. the result is the same. .2) can be seen as a transformation. the point P remains the same. This notion of projection can be extended from lines to vector spaces by the following definition: The projection p of a point b ∈ Rn onto a subspace C is the point in C that is closest to b. the geometric interpretation of equation (3. for unit vectors c. in a fixed reference system. because two interpretations of it are possible. equation (3.

because it also describes what happens to the magnitudes of the vectors: a hypersphere is stretched or compressed into a hyperellipse. in the sense that U T (b − p) = U T (b − U U T b) = U T b − U T b = 0 . This statement is stronger than saying that A maps rowspace(A) into range(A).26 x 2 CHAPTER 3. we have often used geometric intuition to introduce new concepts.2 The Singular Value Decomposition In these notes. Figure 3. This approach is successful when geometry is less cumbersome than algebra. the hyperellipse is an ellipsoid. For this value of p the difference vector b − p is orthogonal to range(U ). The two small boxes are corresponding points. The geometric picture underlying the Singular Value Decomposition is crisp and useful. and we have then translated these into algebraic statements.5) −3 3  2 1 1 transforms the unit circle on the plane into an ellipse embedded in three-dimensional space. which is zero iff x = UT b .1 shows the map b = Ax . which is a quadratic hypersurface that generalizes the two-dimensional notion of ellipse to an arbitrary number of dimensions. THE SINGULAR VALUE DECOMPOSITION b3 x v2 x 1 v1 u3 σ u2 2 σ u1 1 b b 1 b 2 Figure 3. in one dimension it is a pair of points. that is. the rank-2 matrix  √ √  3 1  3 A= √ (3. or when geometric intuition provides a strong guiding element.1: The matrix in equation (3. when p = Ux = UUT b as promised. . In three dimensions. the hyperellipse in question is centered at the origin. For instance. so we will use geometric intuition again. Here is the main intuition: An m × n matrix A of rank r maps the r-dimensional unit hypersphere in rowspace(A) into an rdimensional hyperellipse in range(A). ∆ 3. In all cases.5) maps a circle on the plane into an ellipse in space.

THE SINGULAR VALUE DECOMPOSITION 27 Two diametrically opposite points on the unit circle are mapped into the two endpoints of the major axis of the ellipse. thereby contradicting the fact that uT Av1 is a maximum. we will prove it algebraically by first introducing the existence of the SVD and then using the latter to prove that matrices map hyperspheres into hyperellipses. If this were not the case. and its other entries 1 T are uj Av1 = 0 because Av1 is parallel to u1 and therefore orthogonal. 1 Similarly. y) must achieve a maximum value on S. and two other diametrically opposite points on the unit circle are mapped into the two endpoints of the minor axis of the ellipse. . n) and σ1 ≥ . Let x and y be unit vectors in Rn and Rm . . . x = y = 1} is compact. y | x ∈ Rn . . . σp ) ∈ Rm×n A = U ΣV T .1 and 2.4. vn .2. By theorems 2. where p = min(m. Simple and fundamental as this geometric fact may be. its proof by geometric means is cumbersome. 1 . . their inner product uT Av1 could 1 be increased by rotating u1 towards the direction of Av1 . by noticing that uT Av1 = vT AT u1 1 1 and repeating the argument above. ≥ σp ≥ 0. We can repeat the same construction on A1 and write T U2 A1 V2 = S2 = 4 Actually. we see that v1 is parallel to AT u1 . σ2 0 0T A2 at least at two points: if uT Av1 is a maximum.1 If A is a real m × n matrix then there exist orthogonal matrices U V such that = = u1 v1 ··· ··· um vn ∈ Rm×m ∈ Rn×n U T AV = Σ = diag(σ1 . . Equivalently. . A similar argument shows that the entries after the first in the first row of S1 are zero: the row vector uT A is parallel to vT . so is (−u1 )T A(−v1 ). y ∈ Rm . . . to u2 . so the first entry of U1 AV1 is uT σ1 u1 = σ1 . so that uT Av2 = .2. . . the first column of AV1 is Av1 = σ1 u1 . T In fact. . This result can be generalized to any m × n matrix. The set S = {x. The lines through these two pairs of points on the unit circle are always orthogonal. Proof. and 1 1 therefore orthogonal to v2 . Then T U1 AV1 = S1 = σ1 0 0T A1 . . Instead.3. um . Let u1 . possibly at more than one point 4 . by construction. = uT Avn = 0.4. respectively. and consider the bilinear form z = yT Ax . 1 1 The matrix A1 has one fewer row and column than A. Collect these orthonormal basis vectors into orthogonal matrices U1 and V1 . u1 and v1 can be extended into orthonormal bases for Rm and Rn . 1 x = y =1 It is easy to see that u1 is parallel to the vector Av1 . respectively. so that the scalar function z(x. . Theorem 3. and let σ1 be the corresponding value of z: max yT Ax = uT Av1 = σ1 . v1 be two unit vectors in Rm and Rn respectively where this maximum is achieved. .2.

It only remains to show that the elements on the diagonal of Σ are nonnegative and arranged in nonincreasing order. . . σp ) . ∆ We can now review the geometric picture in figure 3. and induction will do the rest. . y) is positive. v2 so they map into the major and the minor axis. because axis endpoints come in pairs.28 so that 1 0 0 T U2 T T U1 AV1 CHAPTER 3. ˆ The vector x is equal to the unit vector [0 x]T transformed by the orthogonal matrix V1 . so is −x. and imagine transforming point x (the small box at x on the unit circle) into its corresponding point b = Ax (the small box on the ellipse). and Σ = diag(σ1 . . The σi s are maxima of the function z(x. and is therefore itself a unit vector. . There are a few ways to do this. In the process. A2 1 0 0 V2 T σ1 = 0 0 0 σ2 0 This procedure can be repeated until Ak vanishes (zero rows or zero columns) to obtain U T AV = Σ where U T and V are orthogonal matrices obtained by multiplying together all the orthogonal matrices used in the procedure. . In addition. and is therefore orthogonal to v1 . represented by equation (3. we introduce some nomenclature for the three matrices in the SVD. then z(−x. y) which always assumes both positive and negative values on any hyper-sphere: If z(x. u2 on the ellipse are called left singular vectors. Because x and y thus defined belong to subsets (actually sub-spheres) of the unit spheres in Rn and Rm . .1 in light of the singular value decomposition.5). in this order. consider the vectors x = V1 0 ˆ x and y = U1 0 ˆ y . A similar argument shows that y is a unit vector orthogonal to u1 . . n)). To see that σ1 ≥ . Consider the map in figure 3. v2 on the unit circle that map into the major axes of the ellipse. .2): 1. we conclude that σ2 ≤ σ1 . we can observe that the successive maximization problems that yield σ1 . we can premultiply the matrix product in the theorem by U and postmultiply it by V T to obtain A = U ΣV T . Write x in the frame of reference of the two vectors v1 . . We have σ2 = = ˆ x max ˆT A1 ˆ = max y x ˆ ˆ ˆ = y =1 x = y =1 max 0 ˆ y T T U1 AV1 0 0 ˆ x ˆ y = T S1 0 ˆ x ˆ ˆ x = y =1 max yT Ax ≤ σ1 .1. Just pick one way. The corresponding axis unit vectors u1 . it is a linear combination of v2 . . σp are performed on a sequence of sets each of which contains the next. vn . The σi are nonnegative because all these maximizations are performed on unit hyper-spheres. v2 the two right singular vectors of A. x = y =1 xT v1 = yT u1 = 0 To explain the last equality above. This transformation can be achieved in three steps (see figure 3. If we define V = v1 v2 . and if x is on a hyper-sphere. Let us call v1 . . we just need to show that σ2 ≤ σ1 . . . . but order v1 . y) is negative. ≥ σp (where p = min(m. Since matrices U and V are orthogonal. THE SINGULAR VALUE DECOMPOSITION   0T 0T  . To show this. which is the desired result.

2: Decomposition of the mapping in figure 3. Transform ξ into its image on a “straight” version of the final ellipse. y2 axes. so b = Uη where the orthogonal matrix U= collects the left singular vectors of A. THE SINGULAR VALUE DECOMPOSITION x 2 29 y3 x v1 v2 x 1 u3 σu 2 2 σ u1 1 y y 1 y η 2 ξ 2 2 v’2 ξ σ u’ 2 2 v’1 ξ1 η σ u’ 1 1 η1 Figure 3. “Straight” here means that the axes of the ellipse are aligned with the y1 . nonnegative numbers σ1 . 2. u1 u2 u3 .3. the new coordinates ξ of x become ξ = V Tx because V is orthogonal. and maps it to b. Rotate the reference frame in Rm = R3 so that the “straight” ellipse becomes the ellipse in figure 3. The real. The components of η are the signed magnitudes of the projections of b along the unit vectors u1 .2. This rotation brings η along. u2 . σ2 (major axis first). Otherwise.1. the transformed vector η has coordinates η = Σξ where σ1 Σ= 0 0   0 σ2  0 is a diagonal matrix.1. If the lengths of the half-axes of the ellipse are σ1 . the “straight” ellipse has the same shape as the ellipse in figure 3. σ2 are called the singular values of A. u3 that identify the axes of the ellipse and the normal to the plane of the ellipse.1. 3.

The first is in the orientation of the singular vectors. it can be removed by imposing. n). the rightmost m × (n − m) block of Σ is zero. This is a trivial ambiguity. This is the SVD of A. Singular vectors must be flipped in pairs (a left vector and its corresponding right vector) because the singular values are required to be nonnegative. 1 : p). the right subspace and yield. . all of the form I = U IU T where U is any orthogonal matrix of suitable size. the axes of the latter are not defined. = 0 . a valid SVD. Σp is p × p. the singular values of a matrix A are the lengths of the semi-axes of the hyperellipse E defined by E = {Ax : x = 1} . and Vp = V (:. the subspaces identified by the corresponding left and right singular vectors are unique. the bottom (m − n) × n block of Σ is zero. and V is n × n. . Σ has the same shape and size as A. for instance. However. and still obtain a valid SVD. and write A = Up Σp VpT where Up is m × p. the identity matrix has an infinity of SVDs. we can define Up = U (:. 1 : r). SVD. 1 : r). . Moreover. . For instance.” equivalent version of the SVD. so that the last m − n columns of U are multiplied by zero. . If p = min(m. Thus. If desired. if m < n. The sizes of the matrices in the SVD are as follows: U is m × m. provided that the corresponding left singular vector is flipped as well. . . i which is an even smaller. . Σ is m × n. One can flip any right singular vector. . minimal. If the matrix A maps a hypersphere into another hypersphere.30 CHAPTER 3. 1 : p). vn } range(A) = span{u1 . Even in the general case. if σr is the smallest nonzero singular value of A. The singular value decomposition is “almost unique”. while U and V are square. The second source of ambiguity is deeper. . THE SINGULAR VALUE DECOMPOSITION We can concatenate these three transformations to obtain b = U ΣV T x or A = U ΣV T since this construction works for any point x on the unit circle. together with the corresponding left singular vectors. then we have r A = Ur Σr VrT = i=1 σi ui vT . if p − r singular values are zero. ≥ σr > σr+1 = . both the 2-norm and the Frobenius norm m n A F = i=1 j=1 |aij |2 . Σr = Σ(1 : r. we can let Ur = U (:. say. . . that the first nonzero entry of every left singular value be positive. ur } . There are two sources of ambiguity. More generally. Σp = Σ(1 : p. The SVD reveals a great deal about the structure of a matrix. Similarly. 1 : p). This suggests a “small. and Vr = V (:. but any orthonormal basis can be chosen within. Except for these ambiguities. If we define r by σ1 ≥ . and Vp is n × p. then rank(A) = r null(A) = span{vr+1 . the SVD is unique. 1 : r). that is. if m > n. Finally. whenever two or more singular values coincide. and this multiplies the last n − m rows of V .

but this is the best we can do for this problem. this is a rather high residual.6) arising from a real-life application may or may not admit a solution. but rank 2. For instance. see again chapter 2). . Incompatibility and underdeterminacy can occur together: the system admits no solution. An approximate solution. .6) may not exist. yields exactly the same residual as x (check this). where the double bars henceforth refer to the Euclidean norm. Then. For instance. However. that is. x = [0 2 2]. the linear system (3. . it makes more sense to find a vector x that minimizes the norm Ax − b of the residual vector r = Ax − b . and the least-squares solution is not unique. THE PSEUDOINVERSE and A are neatly characterized in terms of the SVD: A 2 F A 2 2 2 = σ1 + . because measurements are unreliable.3 The Pseudoinverse One of the most important applications of the SVD is the solution of linear systems in the least squares sense. but it tries to satisfy all of them as closely as possible. a vector x that satisfies this equation exactly. x cannot exactly satisfy any of the m equations in the system. in the least-squares sense. always exists. and equations are often mutually incompatible because they come from inexact measurements (incompatible linear systems were defined in chapter 2). In other circumstances. In summary. 2 31 = sup x=0 Ax x In the next few sections we introduce fundamental results and applications that testify to the importance of the SVD. not enough measurements are available.and right-hand sides of the equations. Even when m ≤ n the equations can be incompatible. in the least-squares sense). as measured by the sum of the squares of the discrepancies between left. + σp = σ1 . because of errors in the measurements that produce the entries of A. but may fail to be unique.3. any other vector of the form     1 −1 x =  1  + α 1  2 0 is as good as x. 3. or may not be unique. which has norm 2 (admittedly.6) is underdetermined. A linear system of the form Ax = b (3.3. as we learned in chapter 2. an exact solution to the system (3. In these cases. obtained for α = 1. Often more measurements are available than strictly necessary. Thus. the system x1 + x2 x1 + x2 x3 = 1 = 3 = 2 has three unknowns. in the sense that it has fewer independent equations than unknowns (its rank r is less than n. This leads to more equations than unknowns (the number m of rows in A is greater than the number n of columns). A least-squares solution turns out to be x = [1 1 2]T with residual r = Ax − b = [1 − 1 0]. and its first two equations are incompatible: x1 + x2 cannot be equal to both 1 and √ 3.

But orthogonal matrices do not change the norm of vectors they are applied to (theorem 3. 0 . Σy − c . In order to find the solution to this minimization problem. The matrix A† = V Σ† U T 1/σ1 . . then one of them turns out to be shorter than all the others. THE SINGULAR VALUE DECOMPOSITION If there are several least-squares solutions. . yn cm 0 0 .32 CHAPTER 3.2).   . . the shortest vector x that achieves the min Ax − b . 0 0 ··· .1. U (ΣV T x − U T b) (3. One can therefore redefine what it means to “solve” a linear system so that there is always exactly one solution.. 1/σr 0 . its norm x is smallest..        .8) because U is an orthogonal matrix.  . .3.  .      0 is called the pseudoinverse of A. 0 ··· 0     . . U U = I. Proof.  y  r+1  −  cr+1  . that is. all equally good (or bad).  .   .  ..   . This can be written as T U ΣV T x − b . which both proves uniqueness and provides a recipe for the computation of the solution.  .    σr yr   cr        .   . let us spell out the last expression. with y = V T x and c = U T b..   . .7) where       † Σ =     is an n × m diagonal matrix. We want to minimize the norm of the following vector:   σ1 0 ··· 0  y1   c1        0 . .   . x is unique. that is. and is given by ˆ = V Σ† U T b x (3.. Theorem 3. ··· 0  .   . This minimum norm solution is the subject of the following theorem. so that the last expression above equals ΣV T x − U T b or.1 The minimum-norm least squares solution to a linear system Ax = b.  . . . The minimum-norm Least Squares solution to Ax = b is the shortest vector x that minimizes Ax − b that is.

we will not be able to zero these differences. Without any constraints on x. the meaning of a least-squares solution is therefore usually modified. is the norm of Σy − c. x = 0 is the minimum-norm solution to any homogeneous linear system. . the norm of Ax − b when x is the solution vector. by imposing the constraint x =1 (3. . For homogeneous linear systems. we also want the shortest y. .4 Least-Squares Solution of a Homogeneous Linear Systems Theorem 3. . n . yr . .  0− . on the other hand. least-squares solution to the original system is the unique vector ˆ = V y = V Σ+ c = V Σ+ U T b x as promised. which happens to be an exact solution. In conclusion. and each of them contributes a residual |ci | to the solution. and minimum-norm solution forces the other entries of y. so they have no effect on the solution. cm  33 and do not depend on the unknown y.9) . . the minimum-norm. The residual. = yn = 0. Although correct. When written as a function of the vector c. . once more. . since this vector is related to Ax − b by an orthogonal transformation (see equation (3. r σi = 0 for i = r + 1. because x and y are related by an orthogonal transformation. the solution is trivial: the minimum-norm solution to Ax = 0 is x=0. that is. .1 works regardless of the value of the right-hand side vector b. . . by definition). In each of the first r differences. that is.4. . which is therefore unique: minimum residual forces the choice of y1 . we would fall back to x = 0 again. m are nonzero. this is y = Σ+ c . that is. when the system is homogeneous. Thus. .8)). . Since we look for the minimum-norm solution. In summary. Of course it is not necessarily the only one (any vector in the null space of A is also a solution. . LEAST-SQUARES SOLUTION OF A HOMOGENEOUS LINEAR SYSTEMS The last m − r differences are of the form  cr+1  . . but it is obviously the one with the smallest norm. the square of the residual is m m Ax − b 2 = Σy − c 2 = i=r+1 c2 = i (uT b)2 i i=r+1 which is the projection of the right-hand side vector b onto the complement of the range of A. ∆ 3. this solution is not too interesting.  .3. Thus. what is desired is a nonzero vector x that satisfies the system (3. . there is freedom in their choice. When b = 0. . Notice that there is no other choice for y. the last n − r components of y are multiplied by zeros.9) as well as possible. for the shortest vector x. In other words. We therefore set yr+1 = . In many applications. the desired y has the following components: yi yi = ci for i = 1. Thus. there is nothing we can do about those differences: if some or all the ci for i = r + 1.3.

. . . (3. Proof. Unfortunately. THE SINGULAR VALUE DECOMPOSITION on the solution. The unit-norm Least Squares solution to Ax = 0 is the vector x with x = 1 that minimizes Ax that is. this norm is the same as ΣV T x or. let vn−k+1 . ∆ Section 3.1 Let A = U ΣV T be the singular value decomposition of A. = σn . and the unit-norm constraint on y yields equation (3.4. . if k = 1. In any event. . . When A is rank deficient. . that is. Then. 2 2 2 2 σ1 y1 + . . it may often have more than one singular value equal to zero. . on the other case. .12) is equivalent to equation (3. Since V is orthogonal.11) are unit-norm least squares solutions to the homogeneous linear system Ax = 0.1. Furthermore. If k > 1.1. U ΣV T x . x =1 Note: when σn is greater than zero the most common case is k = 1. .4. the resulting constrained minimization problem does not necessarily admit a unique solution. = yn−k = 0. . that is.11). that is by letting y1 = . all vectors of the form x = α1 vn−k+1 + . Theorem 3. that is. . . Since orthogonal matrices do not change the norm of vectors they are applied to (theorem 3. αk = yn . . + αk = 1 (3. x = 1 translates to y = 1. This is obviously achieved by concentrating all the (unit) mass of y where the σs are smallest. they achieve the min Ax . . . since it is very unlikely that different singular values have exactly the same numerical value.2).5 shows a sample use of theorem 3. vn be the k columns of V whose corresponding singular values are equal to the last singular value σn .12) From y = V T x we obtain x = V y = y1 v1 + . We thus look for the unit-norm vector y that minimizes the norm (squared) of Σy.34 CHAPTER 3. .10) (3. The following theorem provides a recipe for finding this solution. . the theorem above shows how to express all solutions as a linear combination of the last k columns of V . x = vn . . and shows that there is in general a whole hypersphere of solutions. so that equation (3. + σn yn . . with y = V T x. Σy . then the minimum-norm solution is unique. let k be the largest integer such that σn−k+1 = . + yn vn . + αk vn with 2 2 α1 + . The reasoning is very similar to that for the previous theorem.10) with α1 = yn−k+1 .

is called the line normal. (3.5.5. The best-fit line minimizes the sum of the squared distances.2 The Best Line Fit ∂ d ∂c 2 Since the third line parameter c does not appear in the constraint (3.16) If we define the centroid p of all the points pi as p= 1 T P 1. (3. SVD LINE FITTING 35 3.15) we must have =0. yi )T and line ax + by − c = 0 is |axi + byi − c|.14) i pi |c| b a Figure 3.3.15). . The distance from the line to the origin is |c| (see figure 3. and the distance between the line n and a point pi is equal to di = |axi + byi − c| = |pT n − c| . and let be the equation of a line. . b)T .1 Fitting a Line to a Set of Points ax + by − c = 0 Let pi = (xi . (3. .3: The distance between point pi = (xi . 2 = min P n − c1 n =1 2 . 3.5. . . m . the best-fit line achieves the min d n =1 In equation (3.13). 1 is a vector of m ones. if we let d = (d1 . Thus. (3. the line does not change. pm )T . dm ) and P = (p1 . we can assume without loss of generality that n = a2 + b2 = 1 . yi )T be a set of m ≥ 2 points on the plane. .15) 3. . at the minimum (3.13) where the unit vector n = (a. If the lefthand side of this equation is multiplied by a nonzero constant.5 SVD Line Fitting The Singular Value Decomposition of a matrix yields a simple method for fitting a line to a set of points on the plane. orthogonal to the line.3). Thus.

we obtain min d n =1 2 = min P n − 1pT n n =1 2 = min Qn n =1 2 . c = pT n . pm )T . . Equivalently. since Qv2 has norm σ2 . . By replacing this expression into equation (3. when n = v2 . to fit a line (a. In fact. where Q = P − 1pT collects the centered coordinates of the m points. the residue is min d = σ2 n =1 and more specifically the distances di are given by d = σ2 u2 where u2 is the left singular vector corresponding to σ2 . We can solve this constrained minimization problem by theorem 3.16) yields ∂ d ∂c 2 CHAPTER 3. m from which we obtain c= that is.36 equation (3.15). To summarize. the SVD 2 Q = U ΣV T = i=1 σi ui vT i yields 2 Qn = Qv2 = i=1 σi ui vT v2 = σ2 u2 i because v1 and v2 are orthonormal vectors.4. . compute the centroid of the points (1 is a vector of m ones): p= 2. b. form the matrix of centered coordinates: Q = P − 1pT 3. we can recall that if n is on a circle. THE SINGULAR VALUE DECOMPOSITION = ∂ nT P T − c1T (P n − 1c) ∂c ∂ = nT P T P n + c2 1T 1 − 2nT P T c1 ∂c = 2 mc − nT P T 1 = 0 1 T T n P 1. compute the SVD of Q: Q = U ΣV T 1 T P 1 m . c) to a set of m points pi collected in the m × 2 matrix P = (p1 . the shortest vector of the form Qn is obtained when n is the right singular vector v2 corresponding to the smaller σ2 of the two singular values of Q.1. and in order to emphasize the geometric meaning of signular values and vectors. proceed as follows: 1. Furthermore.

Here “plus” means the following. % centroid of all the points p = (P’ * one) / m.3. % assemble the three line coefficients into a column vector l = [n . the residue of the fit is The following matlab code implements the line fitting method. A useful exercise is to think how this procedure. if n ˜= 2. Let L be a linear space. 2). % matrix of centered coordinates Q = P . p’ * n]. An affine subspace is a linear subspace plus a point. or something close to it. Hint: minimizing the distance between a point and a subspace is equivalent to maximizing the norm of the projection of the point onto the subspace. just like an arbitrary line is a line through the origin plus a point. 2). 1). b)T = v2 . residue] = linefit(P) % check input matrix sizes [m n] = size(P). can be adapted to fit a set of data points in Rm with an affine subspace of given dimension n. the line normal is the second column of the 2 × 2 matrix V : n = (a. [U Sigma V] = svd(Q). the third coefficient of the line is 37 c = pT n min d = σ2 n =1 6. end one = ones(m. SVD LINE FITTING 4. % the smallest singular value of Q % measures the residual fitting error residue = Sigma(2. The fitting problem (including fitting a line to a set of points) can be cast either as a maximization or a minimization problem. error(’matrix P must be m x 2’). end if m < 2. function [l. 5. . Then an affine space has the form A = p + L = {a | a = p + l and l ∈ L} .one * p’. % the line normal is the second column of V n = V(:.5. error(’Need at least two points’).

38 CHAPTER 3. THE SINGULAR VALUE DECOMPOSITION .

Chapter 4

Function Optimization
There are three main reasons why most problems in robotics, vision, and arguably every other science or endeavor take on the form of optimization problems. One is that the desired goal may not be achievable, and so we try to get as close as possible to it. The second reason is that there may be more ways to achieve the goal, and so we can choose one by assigning a quality to all the solutions and selecting the best one. The third reason is that we may not know how to solve the system of equations f(x) = 0, so instead we minimize the norm f(x) , which is a scalar function of the unknown vector x. We have encountered the first two situations when talking about linear systems. The case in which a linear system admits exactly one exact solution is simple but rare. More often, the system at hand is either incompatible (some say overconstrained) or, at the opposite end, underdetermined. In fact, some problems are both, in a sense. While these problems admit no exact solution, they often admit a multitude of approximate solutions. In addition, many problems lead to nonlinear equations. Consider, for instance, the problem of Structure From Motion (SFM) in computer vision. Nonlinear equations describe how points in the world project onto the images taken by cameras at given positions in space. Structure from motion goes the other way around, and attempts to solve these equations: image points are given, and one wants to determine where the points in the world and the cameras are. Because image points come from noisy measurements, they are not exact, and the resulting system is usually incompatible. SFM is then cast as an optimization problem. On the other hand, the exact system (the one with perfect coefficients) is often close to being underdetermined. For instance, the images may be insufficient to recover a certain shape under a certain motion. Then, an additional criterion must be added to define what a “good” solution is. In these cases, the noisy system admits no exact solutions, but has many approximate ones. The term “optimization” is meant to subsume both minimization and maximization. However, maximizing the scalar function f (x) is the same as minimizing −f (x), so we consider optimization and minimization to be essentially synonyms. Usually, one is after global minima. However, global minima are hard to find, since they involve a universal quantifier: x∗ is a global minimum of f if for every other x we have f (x) ≥ f (x∗ ). Global minization techniques like simulated annealing have been proposed, but their convergence properties depend very strongly on the problem at hand. In this chapter, we consider local minimization: we pick a starting point x0 , and we descend in the landscape of f (x) until we cannot go down any further. The bottom of the valley is a local minimum. Local minimization is appropriate if we know how to pick an x0 that is close to x∗ . This occurs frequently in feedback systems. In these systems, we start at a local (or even a global) minimum. The system then evolves and escapes from the minimum. As soon as this occurs, a control signal is generated to bring the system back to the minimum. Because of this immediate reaction, the old minimum can often be used as a starting point x0 when looking for the new minimum, that is, when computing the required control signal. More formally, we reach the correct minimum x∗ as long as the initial point x0 is in the basin of attraction of x∗ , defined as the largest neighborhood of x∗ in which f (x) is convex. Good references for the discussion in this chapter are Matrix Computations, Practical Optimization, and Numerical Recipes in C, all of which are listed with full citations in section 1.4. 39

40

CHAPTER 4. FUNCTION OPTIMIZATION

4.1 Local Minimization and Steepest Descent
Suppose that we want to find a local minimum for the scalar function f of the vector variable x, starting from an initial point x0 . Picking an appropriate x0 is crucial, but also very problem-dependent. We start from x0 , and we go downhill. At every step of the way, we must make the following decisions: • Whether to stop. • In what direction to proceed. • How long a step to take. In fact, most minimization algorithms have the following structure: k=0 while xk is not a minimum compute step direction pk with pk = 1 compute step size αk xk+1 = xk + αk pk k =k+1 end. Different algorithms differ in how each of these instructions is performed. It is intuitively clear that the choice of the step size αk is important. Too small a step leads to slow convergence, or even to lack of convergence altogether. Too large a step causes overshooting, that is, leaping past the solution. The most disastrous consequence of this is that we may leave the basin of attraction, or that we oscillate back and forth with increasing amplitudes, leading to instability. Even when oscillations decrease, they can slow down convergence considerably. What is less obvious is that the best direction of descent is not necessarily, and in fact is quite rarely, the direction of steepest descent, as we now show. Consider a simple but important case, 1 f (x) = c + aT x + xT Qx 2 (4.1)

where Q is a symmetric, positive definite matrix. Positive definite means that for every nonzero x the quantity xT Qx is positive. In this case, the graph of f (x) − c is a plane aT x plus a paraboloid. Of course, if f were this simple, no descent methods would be necessary. In fact the minimum of f can be found by setting its gradient to zero: ∂f = a + Qx = 0 ∂x so that the minimum x∗ is the solution to the linear system Qx = −a . (4.2)

Since Q is positive definite, it is also invertible (why?), and the solution x∗ is unique. However, understanding the behavior of minimization algorithms in this simple case is crucial in order to establish the convergence properties of these algorithms for more general functions. In fact, all smooth functions can be approximated by paraboloids in a sufficiently small neighborhood of any point. Let us therefore assume that we minimize f as given in equation (4.1), and that at every step we choose the direction of steepest descent. In order to simplify the mathematics, we observe that if we let e(x) = ˜ then we have 1 (x − x∗ )T Q(x − x∗ ) 2 (4.3)

1 e(x) = f (x) − c + x∗ T Qx∗ = f (x) − f (x∗ ) ˜ 2

4.1. LOCAL MINIMIZATION AND STEEPEST DESCENT
so that e and f differ only by a constant. In fact, ˜ e(x) = ˜ 1 T 1 1 1 (x Qx + x∗ T Qx∗ − 2xT Qx∗ ) = xT Qx + aT x + x∗ T Qx∗ = f (x) − c + x∗ T Qx∗ 2 2 2 2 1 1 1 f (x∗ ) = c + aT x∗ + x∗ T Qx∗ = c − x∗ T Qx∗ + x∗ T Qx∗ = c − x∗ T Qx∗ . 2 2 2 The result (4.3), e(x) = f (x) − f (x∗ ) , ˜

41

and from equation (4.2) we obtain

is rather interesting in itself. It says that adding a linear term aT x (and a constant c) to a paraboloid 1 xT Qx merely 2 1 shifts the bottom of the paraboloid, both in position (x∗ rather than 0) and value (c − 2 x∗ T Qx∗ rather than zero). Adding the linear term does not “warp” or “tilt” the shape of the paraboloid in any way. Since e is simpler, we consider that we are minimizing e rather than f . In addition, we can let ˜ ˜ y = x − x∗ , that is, we can shift the origin of the domain to x∗ , and study the function e(y) = 1 T y Qy 2

instead of f or e, without loss of generality. We will transform everything back to f and x once we are done. Of ˜ course, by construction, the new minimum is at y∗ = 0 where e reaches a value of zero: e(y∗ ) = e(0) = 0 .

However, we let our steepest descent algorithm find this minimum by starting from the initial point y0 = x0 − x∗ . At every iteration k, the algorithm chooses the direction of steepest descent, which is in the direction pk = − opposite to the gradient of e evaluated at yk : gk = g(yk ) = ∂e = Qyk . ∂y y=y k gk gk

We select for the algorithm the most favorable step size, that is, the one that takes us from yk to the lowest point in the direction of pk . This can be found by differentiating the function e(yk + αpk ) = 1 (y + αpk )T Q(yk + αpk ) 2 k

with respect to α, and setting the derivative to zero to obtain the optimal step αk . We have ∂e(yk + αpk ) = (yk + αpk )T Qpk ∂α and setting this to zero yields αk = − (Qyk )T pk gT p pT p gT g = − Tk k = gk Tk k = gk Tk k . pT Qpk pk Qpk pk Qpk gk Qgk k (4.4)

5) How much closer did this step bring us to the solution y∗ = 0? In other words. we introduce the following result. e(yk ) gk Q gk gk Qgk This can be rewritten as follows by rearranging terms: e(yk+1 ) = 1− (gT gk )2 k T Q−1 g gT Qg gk k k k e(yk ) (4. .42 Thus. G.6) so if we can bound the expression in parentheses we have a bound on the rate of convergence of steepest descent. relative to the value e(yk ) at the previous step? The answer is. respectively. Luenberger. the largest and smallest singular values of Q.5) we obtain e(yk ) − e(yk+1 ) e(yk ) = yT Qyk − yT Qyk+1 k k+1 yT Qyk k T gT gk gT gk yT Qyk − yk − gTkQg gk Q yk − gTkQg gk k k k k k yT Qyk k = = = Since Q is invertible we have gT gk gT gk 2 gTkQg gT Qyk − gTkQg k k k k k yT Qyk k 2gT gk gT Qyk − (gT gk )2 k k k . yT Qyk gT Qgk k k 2 gT Qgk k gk = Qyk and ⇒ yk = Q−1 gk yT Qyk = gT Q−1 gk k k so that e(yk ) − e(yk+1 ) (gT g )2 = T −1k k T . the basic step of our steepest descent can be written as follows: yk+1 = yk + gk that is. AddisonWesley. n × n matrix. 1973. Introduction to Linear and Nonlinear Programming. To this end. Lemma 4.1 (Kantorovich inequality) Let Q be a positive definite. FUNCTION OPTIMIZATION gT gk k p gT Qgk k k gT gk k g . gT Qgk k k (4. The arguments and proofs below are adapted from D.1. For any vector y there holds 4σ1 σn (yT y)2 ≥ yT Q−1 y yT Qy (σ1 + σn )2 where σ1 and σn are. yk+1 = yk − CHAPTER 4. symmetric. From the definition of e and from equation (4. as we shall now prove. often not much. how much smaller is e(yk+1 ).

the values 1/σj for j = 1.8) then the numerator φ(σ) in (4.7) is a convex combination of points on this curve. yielding the desired result. Because Q is positive definite. Thus. each of the singular values σj can be obtained by letting θj = 1 and all other θi to zero. the method of steepest descent xk+1 = xk − where gk = g(xk ) = gT gk k g gT Qgk k k (4. n are all on the curve 1/σ. .7) must be in the shaded area in figure 4. . Let 43 Q = U ΣU T be the singular value decomposition of the symmetric (hence V = U ) matrix Q. Of course. . For the same vector of coefficients θi . and λ(σ) are on the vertical line corresponding to the value of σ given by (4. If we let z = UT y we have 1/ (yT y)2 (yT U T U y)2 (zT z)2 = T = T −1 T = T Q−1 y yT Qy −1 U T y yT U ΣU T y y y UΣ z Σ z z Σz θi = add up to one. ψ(σ) σ1 ≤σ≤σn λ(σ) σ1 ≤σ≤σn (σ1 + σn − σ)/(σ1 σn ) The minimum is achieved at σ = (σ1 + σn )/2. The denominator ψ(σ) in (4. Thus an appropriate bound is φ(σ) φ(σ) 1/σ ≥ min = min . with Q symmetric and positive definite. ψ(σ).7) is 1/σ. . we can state the main result on the convergence of the method of steepest descent. LOCAL MINIMIZATION AND STEEPEST DESCENT Proof. that is.4.1.9) ∂f = a + Qxk ∂x x=xk . 1/σn ). However. Theorem 4. the values of φ(σ). (4. For any x0 .2 Let 1 f (x) = c + aT x + xT Qx 2 be a quadratic function of x. the values of the denominator ψ(σ) of (4. If we let σ= i=1 n 2 zi z 2 n i=1 θi σi n i=1 θi /σi = φ(σ) ψ(σ) (4.1.7) where the coefficients θi σi . Since 1/σ is a convex function of σ. This area is delimited from above by the straight line that connects point (σ1 .8). there are many ways to choose the coefficients θi to obtain a particular value of σ. by the line with ordinate λ(σ) = (σ1 + σn − σ)/(σ1 σn ) . ∆ Thanks to this lemma. since the smallest of them satisfies σn = min yT Qy > 0 y =1 by the definition of positive definiteness. 1/σ1 ) with point (σn . all its singular values are strictly positive.1.

the closer the fraction (σ1 − σn )/(σ1 + σn ) is to unity. that xk → x∗ . The larger the condition number.12) (f (xk ) − f (x∗ )) . Since the ratio in the last term is smaller than one.3) and (4.44 CHAPTER 4. it follows immediately that f (xk ) − f (x∗ ) → 0 and hence.1: Kantorovich inequality.11) (4. FUNCTION OPTIMIZATION φ. Proof. at every step k there holds f (xk+1 ) − f (x∗ ) ≤ σ1 − σn σ1 + σn 2 (f (xk ) − f (x∗ )) where σ1 and σn are. ∆ The ratio κ(Q) = σ1 /σn is called the condition number of Q. respectively. Furthermore. It is easily seen why this happens in the case . converges to the unique minimum point x∗ = −Q−1 a of f .6) and the Kantorovich inequality we obtain f (xk+1 ) − f (x∗ ) = e(yk+1 ) = = σ1 − σn σ1 + σn 1− 2 (gT gk )2 k T Q−1 g gT Qg gk k k k e(yk ) ≤ 1− 4σ1 σn (σ1 + σn )2 e(yk ) (4. and the slower convergence. By equations (4. the largest and smallest singular value of Q.10) we immediately obtain the expression for steepest descent in terms of f and x.λ λ(σ) ψ(σ) φ(σ) σ1 σ2 σ σn σ Figure 4. since the minimum of f is unique. From the definitions y = x − x∗ and e(y) = 1 T y Qy 2 (4.ψ.

Thus. but can become arbitrarily slow for poorly conditioned Hessians. The minimum of f along that line is tangent to some other. This is defined as the largest value of q for which the lim xk+1 − x∗ xk − x∗ q k→∞ is finite. Theorem 4. In that case. so the higher the order of convergence. in which x is a two-dimensional vector. the distance of xk from x∗ is reduced by the q-th power at every step. This case. which is orthogonal to the isocontour at xk . the analysis above applies once xk gets close enough to the minimum.2. and is called the Hessian of f . and minimization would get there in just one more step. non-quadratic) f . In this case. is consistent with our analysis. but very precarious case. we introduce the notion of the order of convergence. then the first turn would make the new direction point to x∗ .1. when the starting point x0 is at one apex (tip of either axis) of an isocontour ellipse.2 implies that steepest descent has at best a linear order of convergence. To characterize the speed of convergence of different minimization algorithms. In summary. which shows the trajectory xk superimposed on a set of isocontours of f (x). In fact. the greater the condition number κ(Q). the farther away a line orthogonal to an isocontour passes from x∗ . that is. lower isocontour. If β is this limit.1. in which κ(Q) = 1. and the more steps are required for convergence. If isocontours were circles (σ1 = σn ) centered at x∗ . will not pass through x∗ .2: Trajectory of steepest descent. In all other cases. The next step is orthogonal to the latter isocontour (that is.4. the line in the direction pk of steepest descent. In other words. LOCAL MINIMIZATION AND STEEPEST DESCENT 45 x* x0 p0 p1 Figure 4. steepest descent is good for functions that have a well conditioned Hessian near the minimum. as in figure 4. For general (that is. σ1 + σn The more elongated the isocontours. then close to the solution (that is. at every step the steepest descent trajectory is forced to make a ninety-degree turn. There is one good. so that f is well approximated by a paraboloid. the better. the residuals |f (xk ) − f (x∗ )| in the values of the function being minimized converge . for large values of k) we have xk+1 − x∗ ≈ β xk − x∗ q for a minimization method of order q. Q is the matrix of second derivatives of f with respect to x. parallel to the gradient). one iteration will lead to the minimum x∗ . because then σ1 − σn =0. namely.

if we cannot assume convexity. if we can assume that h is convex between α1 and αi . the derivatives of f are close to zero. c) otherwise (a. . Since the gradient of f approaches zero when xk tends to x∗ .46 CHAPTER 4. . Close to the minimum. the arguments xk to f can converge to x∗ even more slowly. but its computation during steepest descent would in general be overkill. In contrast. c) = (u. we can set a = 0. c) = (b. so |f (xk ) − f (xk−1 )| may be very small but xk − xk−1 may still be relatively large. b. FUNCTION OPTIMIZATION linearly. Using line search to find αk guarantees that a minimum in the direction pk is actually reached even when the parabolic approximation is inadequate. b. b) end otherwise u = (b + c)/2 if f (u) > f (b) (a. b. Another is to check whether xk is significantly different from xk−1 . corresponding through equation (4. b. rather than in that of f (x∗ ). A point c that is on the opposite side of the minimum with respect to a can be found by increasing α through values α1 = a. the steepest descent algorithm can be stopped when xk − xk−1 < where the positive constant is provided by the user. Then. To complete the steepest descent algorithm we need to specify how to check whether a minimum has been reached. we may find the wrong minimum.13) to the starting point xk . We used Q because it was available. usually one is interested in the value of x∗ . The only difficulty here is to find c. . we used the Hessian Q in order to compute the optimal step size α (see equation (4. but there is no general-purpose fix to this problem. Line search first determines two points a.13) be the scalar function of one variable that is obtained by restricting the function f to the line through the current point xk and in the direction of pk . the check on xk is more stringent. but it is increasing between αi−1 and αi = c. b = (a + c)/2. Of course. Thus. until h(αi ) is greater than h(αi−1 ). c) = (a. In fact. In fact. we can set c = αi . Let h(α) = f (xk + αpk ) (4. Line search now proceeds by shrinking the bracketing triple (a. and therefore preferable in most cases. In fact. so the minimum must be somewhere between a and c. c) until c − a is smaller than the desired accuracy in determining αk .4)). u) otherwise (a. In fact. α2 . In summary. the derivative of h at a is negative. the Hessian of f (x) requires computing n second derivatives if x is an n-dimensional 2 vector. and a line search in the direction of pk can be used to determine the step size αk . Here is how line search works. In our analysis of steepest descent. c that bracket the desired minimum αk . u. only gradient information is necessary to find pk . c) end end. b. c) = (a. One criterion is to check whether the value of f (xk ) has significantly decreased from f (xk−1 ). u. b. in the sense that a ≤ αk ≤ c. so the function is initially decreasing. say. b. . Shrinking works as follows: if b − a > c − b u = (a + b)/2 if f (u) > f (b) (a. and then picks a point between a and c.

the gradient gk is different from a as used in equation (4. . the new value x1 will be different from x∗ . and obtain. Let 1 f (xk + ∆x) ≈ f (xk ) + gT ∆x + ∆xT Qk ∆x k 2 be the first terms of the Taylor series expansion of f about the current point xk . A minimization algorithm in which the step direction pk and size αk are defined in this manner is called Newton’s method. so line search will find the minimum in a number of steps that is logarithmic in the desired accuracy. and therefore the minimum is somewhere between a and c. In other words.15) whose solution ∆xk = αk pk yields at the same time the step direction pk = ∆xk / ∆xk and the step size αk = ∆xk . NEWTON’S METHOD 47 It is easy to see that in each case the bracketing triple (a.2. and suppose that at the minimum x∗ the Hessian Q(x∗ ) is nonsingular. . .4. Newton’s method is quadratic. we can set the derivatives of f (xk + ∆x) with respect to ∆x to zero. a and Q are the coefficients of the Taylor expansion of f around point x = 0. .15)). In addition. but convergence can be expected to be faster. gradient and Hessian are constantly reevaluated in Newton’s method. which we now summarize. In fact. 4. . In fact. Notice that even when f is a paraboloid. This is the idea of Newton’s method. ∂2f ∂x2 n ∂2f ∂xn ∂x1 are the gradient and Hessian of f evaluated at the current point xk .14) is valid. when f (x) is not exactly a paraboloid. that is. at every step the interval (a. Then y(x∗ ) = x∗ because g(x∗ ) = 0. In fact. the analysis in the previous section suggests a straight-forward fix to the slow convergence of steepest descent. The greater speed of Newton’s method over steepest descent is borne out by analysis: while steepest descent has a linear order of convergence. the linear system Qk ∆x = −gk . when αk = 0. where gk = g(xk ) = and Qk = Q(xk ) =   ∂2f = T ∂x∂x x=xk  ∂f ∂x x=xk ∂2f ∂x2 1 (4. b. c) shrinks to 3/4 of its previous size. and xk+1 − x∗ = y(xk ) − x∗ = y(xk ) − y(x∗ ) . The corresponding pk is termed the Newton direction. .2 Newton’s Method If a function can be well approximated by a paraboloid in the region in which minimization is performed. let y(x) = x − Q(x)−1 g(x) be the place reached by a Newton step starting at x (see equation (4.14) ··· ··· ∂2f ∂x1 ∂xn     x= xk . Of course. Consequently. while gk and Qk are the coefficients of the Taylor expansion of f around the current point xk . c) preserves the property that f (b) ≤ f (a) and f (b) ≤ f (c). The direction is of course undefined once the algorithm has reached a minimum.2) tells us how to jump in one step from the starting point x0 to the minimum x∗ .15) is the Newton step.2). (4. and the step defined by equation (4. iterations are needed. equation (4. analogously to equation (4.1). To the extent that approximation (4.

1). and then jumps at every iteration to the lowest point of the current approximation. However. However. conjugate gradients is a direct method for the quadratic case. This paradox. For a quadratic function. for a total of n + n derivatives. a quarter of a million pixels. is only apparent. because it only uses gradient information. is motivated by the desire to accelerate convergence with respect to the steepest descent method. 000 unknown motion field values.48 From the mean-value theorem. It effectively solves the linear system (4. The conjugate gradients method described in these notes is the so-called Polak-Ribi` re variation. Prima facie. it will be developed for the simple case of minimizing a quadratic function with positivedefinite and known Hessian. We know that in this case minimizing f (x) is equivalent to solving the linear system (4. although the version presented here will only work if the matrix of the system is symmetric and positive definite. system (4. each requiring one gradient computation and one line search. where n is the number of components of x. It will be introe duced in three steps. so that x the first term in the right-hand side above vanishes. and a line search in the direction pk to find the step size. In 2 addition. as in equation (4. 4. . These high-dimensional problems arise typically from the discretization of partial differential equations. conjugate gradients saves the day. the convergence rate of Newton’s method is of order at least two. At every pixel in the image. similarly to Newton’s method. but it does so by a sequence of n one-dimensional minimizations. This quadratic function f (x) was introduced in equation (4. Thus. or. say. we have xk+1 − x∗ = y(xk ) − y(x∗ ) ≤ CHAPTER 4. FUNCTION OPTIMIZATION ∂y 1 ∂2y (xk − x∗ ) + xk − x∗ T ∂x x=x∗ 2 ∂x∂xT x=x ˆ 2 where ˆ is some point on the line between x∗ and xk . Storing and inverting a 500. For very large problems. 000 Hessian is out of the question. while Newton’s method reaches the minimum in one step. All it can do is to go downhill. the Hessian must be inverted.15) must be solved. steepest descent takes many steps to converge. First. Specifically. Say for instance that we want to compute the motion of points in an image as a consequence of camera motion. The bottom line is that fast convergence requires work that is equivalent to evaluating the Hessian of f . this motion. Conjugate gradients works by taking n steps for each of the steps in Newton’s method. storing and manipulating a Hessian can be prohibitive. the first derivatives of y at x∗ are zero. in which the dimension n of x is thousands or more. discussed in the next section.2) is never constructed explicitly. In cases like these. The method of conjugate gradients. Since y(x∗ ) = x∗ . the work done by conjugate gradients is equivalent to that done by Newton’s method. because it requires both the gradient gk and the Hessian Qk to be evaluated. called the motion field.1). Thus. Overall. In contrast. This means that the number of iterations is fixed.2). if an image has. This is very important in cases where x has thousands or even millions of components. and xk+1 − x∗ ≤ c xk − x∗ 2 where c depends on third-order derivatives of f near x∗ .2) of Newton’s method. and the matrix Q is never stored. Steepest descent locally approximates the function with planes. at least. but it only requires gradient information. steepest descent requires the gradient gk for selecting the step direction pk . 000 × 500. Newton’s method approximates f with paraboloids. system (4.3 Conjugate Gradients Newton’s method converges faster (quadratically) than steepest descent (linear convergence rate) because it uses more information about the function f being minimized. but without paying the storage cost of Newton’s method. the method of conjugate gradients discussed in this section seems to violate this principle: it achieves fast. there are n = 500. superlinear convergence. Partial differential equations relate image intensities over space and time to the motion of the underlying image features. is represented by a vector whose magnitude and direction describe the velocity of the image feature at that pixel. however. the method converges to the solution in n steps. Rather than an iterative method. conjugate gradients for the quadratic case can also be seen as an alternative method for solving a linear system. this single iteration in Newton’s method is more expensive. Because of the equivalence with a linear system.

We rely instead on the intuition that conjugate gradients solves system (4.16) where Q is a symmetric. on the other hand. so that the new direction p1 is orthogonal to the previous direction p0 . The next step is in the direction of the gradient. . the minimum x∗ is the solution to the linear system Qx = −a . This occurs because the Hessian of f is no longer a constant. the assumption that the Hessian Q in expression (4. but only the quantity gk = Qxk + a that is. positive definite matrix. if f is nearly quadratic in a large neighborhood of the minimum. since p0 . .4. In this case. This combination of fast convergence. When the ellipsoids are spheres. (4.17) We know how to solve such a system. the convergence rate of conjugate gradients is superlinear. CONJUGATE GRADIENTS 49 Second. It is easy to see that in that case n steps will take us to x∗ . somewhere between Newton’s method and steepest descent. In summary. We now consider an alternative solution method that does not need Q. The first step takes from x0 to x1 .3. each line minimization is independent of the others: The first step yields the minimum in the space spanned by p0 . In spite of this.3. We will see that the conjugate gradients method requires n gradient evaluations and n line searches in lieu of each n × n matrix inversion in Newton’s method. and we will omit this proof.1) is known will be removed. As discussed above. pn−1 span the whole space. will lead to the minimum for spherical isosurfaces. the gradient of f (x). After n steps we must be done.1 The Quadratic Case 1 f (x) = c + aT x + xT Qx 2 Suppose that we want to minimize the quadratic function (4. that is. called isosurfaces. the second step then yields the minimum in the space spanned by p0 and p1 . As we saw in our discussion of steepest descent. in which the variable x in f (x) is a three-dimensional vector. and that the quadratic approximation becomes more and more valid as the algorithm converges to the minimum. and low computational cost per iteration explains the popularity of conjugate gradients methods for the optimization of functions of a large number of variables. Consider the case n = 3. Formal proofs can be found in Elijah Polak. any set of orthogonal directions. and so forth. 1997. as it was in the quadratic case. evaluated at n different points x1 . Given an arbitrary set of ellipsoidal isosurfaces. and the line between x0 and x1 is tangent to an isosurface at x1 . all the methods we have seen so far involve explicit manipulation of the matrix Q. xn . Springer. and each of the steps is simple. This would then take us to x∗ right away. and x has n components. 4. Optimization — Algorithms and consistent approximations. The arguments offered below appeal to intuition. convergence is fast. positive definite matrix Q. Finding tight bounds for the convergence rate of conjugate gradients is hard. In fact. and the cost of finding the minimum depends on the desired accuracy. centered at the minimum x∗ . . to slow convergence. that is. As a consequence. Suppose however that we cannot afford to compute this special direction p1 orthogonal to p0 . we noticed that for poorly conditioned Hessians orthogonal directions lead to many small steps. but that we can only compute some direction p1 orthogonal to p0 (there is an n − 1-dimensional space of such directions!). this works much better. then we can write 1 1 T x Qx = yT y 2 2 . as it requires close to the n steps that are necessary in the quadratic case. a certain property that holds in the quadratic case is now valid only approximately.2). modest storage requirements. If the function f starts to behave like a quadratic function early. the method is no longer direct. NY. . but iterative. . this is the main reason for using conjugate gradients. Then the quadratic function f (x) is constant over ellipsoids. How can we start from a point x0 on one of these ellipsoids and reach x∗ by a finite sequence of one-dimensional searches? In connection with steepest descent. since isosurfaces are spheres. However. the conjugate gradients method will be extended to general functions f (x). Third. . there is a one-to-one mapping with a spherical system: if Q = U ΣU T is the SVD of the symmetric. . with a line search in each direction.

If two directions qi and qj are orthogonal in the spherical context. if qT qj = 0 . Here is the procedure. . . In fact. . or Q-orthogonality: if equation (4. which also incorporates the steps from x0 to xn : g0 = g(x0 ) p0 = −g0 for k = 0 . pn−1 that are mutually conjugate. then pi and pj are said to be Q-conjugate or Q-orthogonal to each other.18) in the algorithm. Bureau National Standards. n − 1 αk = arg minα≥0 f (xk + αpk ) xk+1 = xk + αk pk gk+1 = g(xk+1 ) gT Qp k+1 γ k = pT Q p k k k pk+1 = −gk+1 + γk pk end where gk = g(xk ) = is the gradient of f at xk . . i what does this translate into in terms of the directions pi and pj for the ellipsoidal problem? We have qi. if we can find n directions p0 . . Details can be found in Polak’s book mentioned earlier. First. pp. in the next subsection. k are also conjugate. . we cannot use the transformation (4. Then.19) This condition is called Q-conjugacy. section B. 1952). In summary. k k+1 It is somewhat more cumbersome to show that pi and pk+1 for i = 0. We do this in two steps. J. there must be a condition for the original problem (in terms of Q) that is equivalent to orthogonality for the spherical problem.19) holds. Of course. we find conjugate directions whose definitions do involve Q. due to Hestenes and Stiefel (Methods of conjugate gradients for solving linear systems. pT Qpk+1 k = pT Q(−gk+1 + γk pk ) k = −pT Qgk+1 + k gT Qpk T k+1 p Qpk pT Qpk k k ∂f ∂x x=xk = −pT Qgk+1 + gT Qpk = 0 . (4. Res. This can be done by induction. because Σ and especially U T are too large. . we reach the minimum in at most n steps. we rewrite these expressions without Q. FUNCTION OPTIMIZATION y = Σ1/2 U T x . 409-436. and if we do line minimization along each direction pk . Vol 49. The proof is based on the observation that the vectors pk are found by a generalization of Gram-Schmidt (theorem 2.j . i (4. . . .18) Consequently.j becomes pT U Σ1/2 Σ1/2 U T pj = 0 i or pT Qpj = 0 . so that orthogonality for qi.4. So now we need to find a method for generating n conjugate directions without using either Q or its SVD.2) to produce conjugate rather than orthogonal vectors. that is. .50 where CHAPTER 4.j = Σ1/2 U T pi. It is simple to see that pk and pk+1 are conjugate. We will henceforth simply say “conjugate” for brevity. .

2 Removing the Hessian The algorithm shown in the previous subsection is a correct conjugate gradients algorithm. In fact.3. However.3. or αk Qpk = gk+1 − gk . now is a function of xk .3. To this end. pT Qpk pT αk Qpk pT (gk+1 − gk ) k k k and Q has disappeared. pT gk = 0 . without ever constructing the Hessian explicitly. the Hessian.4. We now show that γk can be rewritten in terms of the gradient values gk and gk+1 only. In fact. We can therefore write γk = gT αk Qpk gT (gk+1 − gk ) gT Qpk k+1 = k+1 = k+1 . . we notice that gk+1 = gk + αk Qpk .3 Extension to General Functions We now know how to minimize the quadratic function (4. g(x) = a + Qx so that gk+1 = g(xk+1 ) = g(xk + αk pk ) = a + Q(xk + αk pk ) = gk + αk Qpk . k−1 Then. However. n iterations will not suffice. the denominator of γk becomes pT (gk+1 − gk ) = −pT gk = (gk − γk−1 pk−1 )T gk = gT gk . we obtain the Polak-Ribi` re formula e γk = gT (gk+1 − gk ) k+1 . since pk and pk+1 are associated to different Hessians. and a few cycles of n iterations each will achieve convergence. Strictly speaking. it is computationally inadequate because the expression for γk contains the Hessian Q. When the function f (x) is arbitrary. gT gk k 4. we then lose conjugacy. This expression for γk can be further simplified by noticing that pT gk+1 = 0 k because the line along pk is tangent to an isosurface at xk+1 . which was constant for the quadratic case. Similarly. which is too large. while the gradient gk+1 is orthogonal to the isosurface at xk+1 . CONJUGATE GRADIENTS 51 4. the same algorithm can be used.16) in n steps. k k k In conclusion. However. the quadratic approximation becomes more and more valid. as the algorithm approaches the minimum x∗ .

FUNCTION OPTIMIZATION .52 CHAPTER 4.

however. the singular value decomposition A = U ΣV of A transforms the domain of the transformation via the matrix V T and its range via the matrix U T so that the transformed system is diagonal. This is a fundamental transformation to use whenever the domain and the range of A are separate spaces. if V is an orthogonal matrix and we define y = V Tx . Ideally. ˙ ˙ then the definition of x forces us to transform x by V T as well: dx dV Tx ˙ =VT = V Tx . domain and range are intimately related to one another even independently of the transformation A. we would like to find an orthogonal matrix S and a diagonal matrix Λ such that A = SΛS T 53 (5.1) . the fact that A is square is not a coincidence. T and where Σ is diagonal. pre-existing relation between x and x. the equation b = U ΣV T x can be written as follows U T b = ΣV T x . Often. and one cannot change coordinates for x ˙ without also changing those for x accordingly. while we need a single transformation. x is assumed to be a function ˙ of some real scalar variable t (often time). dt dt In brief.Chapter 5 Eigenvalues and Eigenvectors Given a linear transformation b = Ax . and x is the derivative of x with respect to t: ˙ x= dx . c = Σy where y = V Tx and c = UT b . of the form ˙ x = Ax where A is n × n. In fact. the SVD does nothing useful for systems of linear differential equations. that is. because it diagonalizes A by two different transformations. For this equation. one for the domain and one for the range. In fact. The most important example is perhaps that of a system of linear differential equations. dt ˙ In other words. In fact. there is an intimate.

rather than orthogonal. in order to diagonalize a square matrix A from a system of linear differential equations we generally look for a decomposition of A of the form A = QΛQ−1 (5.2) . Finally. they can be even non-real.1.54 so that if we define CHAPTER 5. Unfortunately. because now we are imposing stronger constraints on the terms of the decomposition. in the sense that they do not alter the vectors’ norms. as well as the inner products between vectors. For complex vectors. scalar differential equations. This will only occur for very special matrices. • the elements on the diagonal of Λ are allowed to be negative. In some cases. rather than real. EIGENVALUES AND EIGENVECTORS y = ST x we can write the equivalent but diagonal differential system ˙ y = Λy . in fact. Unitary matrices merely rotate or flip vectors.1) possible. so unitary generalizes orthogonal for complex matrices. the norm squared is defined as x and if S is unitary we have Sx 2 2 = xH x . for complex matrices we generalize the notion of transpose by introducing the Hermitian operator: The matrix QH (pronounced “Q Hermitian”) is defined to be the complex conjugate of the transpose of Q. A matrix S is said to be unitary if S H S = SS H = I . and the constraint U = V implies that the vectors vi on the circle that map into the axes σi ui of the ellipse are parallel to the axes themselves. we weaken the constraints in several ways: • the elements of S and Λ are allowed to be complex. which can be solved separately.1) is not always possible. To distinguish invertible from orthogonal matrices we use the symbol Q for invertible and S for orthogonal. because it is a system of n independent. We will see all of this in greater detail soon. a nonunitary transformation Q can change the norms of vectors. it will be possible to diagonalize A by orthogonal transformations S and S T . now the circle and the ellipse live in the same space. If we refer back to figure 3. conjugate transposition becomes simply transposition. It is like doing an SVD but with the additional constraint U = V . if x1 and x2 are mutually orthogonal. In summary. If Q happens to be real. writing A in the form (5. In order to make a decomposition like (5. A matrix that is equal to its Hermitian is called a Hermitian matrix. This is now much easier to handle. This stands to reason. in the sense that xH x2 = 0 . 1 then Sx1 and Sx2 are orthogonal as well: xH S H Sx2 = xH x2 = 0 . 1 1 In contrast. Furthermore. • S is required to be only invertible. 2 = xH S H Sx = xH x = x . The solutions can then be recombined through x = Sy . so the Hermitian is a generalization of the transpose.

In contrast. and the diagonal entries of Λ are called eigenvalues. the columns of qi of Q and the diagonal entries λi of Λ are solutions of the eigenvalue/eigenvector equation Ax = λx .5 −1 −1. Figure 5. The columns of Q are called eigenvectors. (5. . As we know.5 0 −0. λn ) . Each point on the unit circle is transformed into some point on the ellipse. diagonal or not. an invertible linear transformation transforms the unit circle into an ellipse. (5. Whenever two matrices A and B. and the transformation of B into A (and vice versa) is called a similarity transformation.5 1 0. .5) on a sample of points on the unit circle. 2 1. The equation A = QΛQ−1 can be rewritten as follows: AQ = QΛ or separately for every column of Q as follows: Aqi = λi qi where Q= q1 ··· qn and Λ = diag(λ1 . . For some special matrices.5) A= 0 2 . and Λ is diagonal.55 where Q and Λ are complex.1: Effect of the transformation (5.3) Thus. Q is invertible. The dashed lines are vectors that do not change direction under the transformation. are related by A = QBQ−1 .5 −2 −2 −1 0 1 2 Figure 5. they are said to be similar to each other. That real eigenvectors and eigenvalues do not always exist can be clarified by considering the eigenvalue problem from a geometrical point of view in the n = 2 case.1 shows the effect of the transformation represented by the matrix √ 2/3 4/ 3 (5. this may specialize to A = SΛS H with unitary S. we have derived this equation from the requirement of diagonalizing a matrix by a similarity transformation.4) which is how eigenvalues and eigenvectors are usually introduced. .

Notice that they do not correspond to the axes of the ellipse. • det(BC) = det(B) det(C).7). for system (5. and the matrix A cannot be diagonalized. so x can be assumed to be a unit vector without loss of generality. Volumes have been written about the properties of the determinant. In other words.6) This is a homogeneous. defined as the (n − 1) × (n − 1) matrix obtained by removing row i and column j from B. it follows. . (5. Each rotation yields another matrix A. equation (5. pulling the solid lines along. in the sense that det(A − λI) = (−1)n (λ − λ1 ) · . has n complex solutions. .4) is homogeneous in x. and that the coefficient of λn is 1.7) From the definition of determinant.56 CHAPTER 5.1. We will see that the answer depends on the matrix A.3)). because of the following results. there are n real. . Equation (5. In figure 5. In fact. but not enough eigenvectors can be found. the eigenvalues and their eigenvectors exist. · (λ − λn ) where some of the λi may coincide. and the problem is best treated as an algebraic one. Therefore. Thus. and that they are not orthogonal. Notice that there are many transformations that map the unit circle into the same ellipse. as evident from figure 5. The case of exactly n distinct eigenvalues is of particular interest. . (5. but the point-to-point transformation is not. Bij is the algebraic complement of entry bij . The geometric intuition is hidden. Matrices represent point-to-point transformations. the two eigenvectors are shown as dashed lines. det(B) = b11 n i+1 bi1 i=1 (−1) det(Bi1 ) if B is 1 × 1 otherwise is zero. and that a rather diverse array of situations may arise. orthonormal eigenvectors. we have to give up the idea of diagonalizing A.1 Computing Eigenvalues and Eigenvectors Algebraically Let us rewrite the eigenvalue equation Ax = λx as follows: (A − λI)x = 0 .7) is a polynomial of degree n in λ.1 can be rotated. bi ··· bj ··· bn ) = − det( b1 ··· bj ··· bi ··· bn ). EIGENVALUES AND EIGENVECTORS for a sample of points on the unit circle. the curve-to-curve transformation from circle to ellipse is unique. perhaps all real. which is called the characteristic equation of A. Given that the directions of the input vectors are generally changed by the transformation A. In particularly good cases. but the resulting ellipse is unchanged. that the left-hand side of equation (5. In some cases. it is sufficient to recall the following properties from linear algebra: • det(B) = det(B T ). q2 that are mapped into themselves by the original transformation A (see equation (5. which admits nontrivial solutions iff the matrix A − λI is rankdeficient. A square matrix B is rank-deficient iff its determinant. the circle in figure 5. This turns out to be good enough for solving linear differential systems. . In bad cases. an n × n matrix has at most n distinct eigenvalues. In this expression. In other cases. square system of equations. just as triangularization was sufficient for solving linear algebraic systems. we need det(A − λI) = 0 . .6) to admit nontrivial solutions. but they are complex. all eigenvalues exist. . it is not obvious whether the eigenvalue problem admits a solution at all. and we can only triangularize it. bn are linearly dependent. In other words. For our purposes. • det( b1 • det( b1 ··· ··· bn ) = 0 iff b1 . 5. by very simple induction.1. The eigenvalue problem amounts to finding axes q1 .

. Let λ and x be an eigenvalue of A and a corresponding eigenvector: Ax = λx . from c1 x1 + . the differences in parentheses are all nonzero. . . . + ck−1 xk−1 = 0 . so that c2 x2 + .1. . λk . ∆ (5. . .8) For Hermitian matrices (and therefore for real symmetric matrices as well). . we obtain xH Ax xH Ax which implies that = = λxH x λ∗ xH x (5. xk are eigenvectors corresponding to eigenvalues λ1 . We can repeat this procedure until only one term remains. the equation c1 x1 + . . . If we multiply equation (5. However. . . . + ck−1 (λk−1 − λk )xk−1 = 0 .10) λxH x = λ∗ xH x . . . . + ck λk xk = 0 and subtracting this equation from equation (5. and this forces c1 = 0. By taking the Hermitian we obtain xH AH = λ∗ xH . Thus.10) from the right by x.9) from the left by xH and equation (5.2 A Hermitian matrix has real eigenvalues. .8) we have c1 (λ1 − λk )x1 + . . . that is. . . . Proof. we have c1 λ1 x1 + . . . + ck λk xk = 0 . Proof. xk corresponding to distinct eigenvalues λ1 . λk are linearly independent. Theorem 5.1. . Since all λi are distinct. and we can replace each xi by xi = (λi − λk )xi . .5. + ck xk = 0 implies that c1 = . . . . . COMPUTING EIGENVALUES AND EIGENVECTORS ALGEBRAICALLY 57 Theorem 5. . . . the last equation can be rewritten as follows: xH A = λ∗ xH . = ck = 0. xk are linearly independent. we have reduced the summation to one containing k − 1 terms. that the vectors x1 . In summary. + ck xk = 0 This entire argument can be repeated for the last equation. (5. .9) Since A = AH . . . Suppose that c1 x1 + . . therefore forcing c2 = 0.1. + ck xk = 0 where the xi are eigenvectors of a matrix A. . + ck xk = 0 we also have c1 λk x1 + . + ck Axk = 0 and because x1 . = ck = 0. . and so forth. A matrix A is Hermitian iff A = AH . By multiplying by A we obtain c1 Ax1 + . . the situation is even better. We need to show that then c1 = .1 Eigenvectors x1 . . which is still an eigenvector of A: c1 x1 + . .

Since the eigenvalue equation Ax = λx is homogeneous in x. Notice that the converse is not true: a matrix can have coincident eigenvalues and still admit n independent.4 Eigenvectors corresponding to distinct eigenvalues of a Hermitian matrix are mutually orthogonal. For instance. any square matrix with n distinct eigenvalues can be diagonalized by a similarity transformation.2 show that when some eigenvalues coincide. EIGENVALUES AND EIGENVECTORS Since x is an eigenvector. and any square Hermitian matrix with n distinct eigenvalues can be diagonalized by a unitary similarity transformation. Proof. however. the eigenvectors can be made to be orthonormal. ∆ Since the two eigenvalues are distinct. ∆ Corollary 5.4. eigenvectors. From theorem 5. λ − µ is nonzero. Consequently.2 µ = µ∗ . The examples in section 5.58 CHAPTER 5.1.1.3 A real and symmetric matrix has real eigenvalues.1. A real and symmetric matrix is Hermitian. Proof. the vector x can be normalized without violating the equation. ∆ Theorem 5. Theorem 5. we obtain yH Ax yH Ax which implies or = = λyH x µyH x . and we must have yH x = 0. First.1. the n × n identity matrix has n equal eigenvalues but n orthonormal eigenvectors (which can be chosen in infinitely many ways).1. the eigenvectors of an n × n Hermitian matrix with n distinct eigenvalues are all mutually orthogonal. . rather diverse situations can arise concerning the eigenvectors.1. and let x and y be corresponding eigenvectors: Ax Ay = = λx µy ⇒ yH A = µyH because A = AH and from theorem 5. If we multiply these two equations by yH from the left and x from the right.6 The determinant of a triangular matrix is the product of the elements on its diagonal.5 An n × n Hermitian matrix with n distinct eigenvalues admits n orthonormal eigenvectors. λyH x = µyH x (λ − µ)yH x = 0 . respectively. Let λ and µ be two distinct eigenvalues of A. so that we have λ = λ∗ as promised. the scalar xH x is nonzero. we point out a simple but fundamental fact about the eigenvalues of a triangular matrix. Proof. Corollary 5. and even orthonormal. ∆ In summary.

If A is triangular. Matrices with n orthonormal eigenvectors are called normal.1. eλn t   c . so is B = A − λI. and from the previous theorem we obtain det(A − λI) = (a11 − λ) · . ..2. ∆ Corollary 5. .11) 0 1 has eigenvalues 2 and 1 and eigenvectors s1 = 1 0 s2 = 0 1 . . ann . Without loss of generality. This follows immediately from the definition of determinant. . the only possibly nonzero bi1 of the matrix B is b11 . so this result holds for diagonal matrices as well. . which because of the properties above has the same eigenvalues. . ∆ Note that diagonal matrices are triangular.7 The eigenvalues of a triangular matrix are the elements on its diagonal. By repeating the argument for B11 and so forth until we are left with a single scalar. GOOD AND BAD MATRICES 59 Proof.5. Then.2 Good and Bad Matrices Solving differential equations becomes much easier when matrices have a full set of orthonormal eigenvectors. . Proof. . we can assume a triangular matrix B to be upper-triangular. we obtain det(B) = b11 · . the matrix 2 0 A= (5. for otherwise we can repeat the argument for the transpose. · (ann − λ) which is equal to zero for λ = a11 . and the summation in the definition of determinant given above reduces to a single term: det(B) = b11 b11 det(B11 ) if B is 1 × 1 otherwise . 5. Normal matrices are good news. The eigenvalues of a matrix A are the solutions of the equation det(A − λI) = 0 . For instance. . because then the n × n system of differential equations ˙ x = Ax has solution n  x(t) = i=1  ci si eλi t = S  eλ1 t . · bnn .

that is. The simplest example of a defective matrix is 0 1 0 0 which has double eigenvalue 0 and only eigenvector [1 0]T . eλn t Computationally this is more expensive. repeated eigenvalues are not a sufficient condition for defectiveness. but independent). More compactly. things can be worse yet.1. However. not all matrices are as good as these.. there may still be a complete set of n eigenvectors. In fact. q1 and q2 are not orthogonal to each other. is that it have repeated eigenvalues. . to have fewer than n eigenvectors.. First. An example of such a matrix is 2 0 −1 1 √ q2 =   H  S x(0) . Its two eigenvectors are     √ 1 1 2  1 q1 =  0  . q2 = 2 0 0 corresponding to eigenvalues 0 and 2 in this order. . has fewer than n eigenvectors. It belongs to the scum of all matrices:   0 2 −1 A= 0 2 1  . we have seen (theorem 5. x(t) = Q   Q x(0) . . because a computation of a Hermitian is replaced by a matrix inversion. while 3 1 0 3 has double eigenvalue 3 and only eigenvector [1 0]T .1) that a matrix with distinct eigenvalues (zero or nonzero does not matter) has a full set of eigenvectors (perhaps nonorthogonal. but they may not be orthonormal. and there is no q3 .   x(t) = S  eλ1 t . as we now show. and 2. A has to have a repeated eigenvalue if it is to be defective. and the solution becomes  λt  e 1   −1 . eλn t Fortunately these matrices occur frequently in practice. λi are the eigenvalues. which has eigenvalues 2 and 1 and nonorthogonal eigenvectors q1 = 1 0 2 2 1 1 . A necessary condition for an n × n matrix to be defective. because the unitary matrix S is replaced by an invertible matrix Q. 0 0 2 Its eigenvalues are 0. Furthermore. and the eigenvectors it has are not orthogonal. EIGENVALUES AND EIGENVECTORS where S = [s1 · · · sn ] are the eigenvectors. as the identity matrix proves. because the matrix is singular. This is conceptually only a slight problem. However. How bad can a matrix be? Here is a matrix that is singular.60 CHAPTER 5. so zero eigenvalues are not the problem. However. repeated twice. and a full set of eigenvectors may fail to exist. and the vector c of constants ci is c = S H x(0) .

5. these matrices can be triangularized by similarity transformations.  . we have 1    T   T  x s1 s1 x    . T y = λy where y = Q−1 x . It may be less obvious that a complex vector x can be transformed into a real vector parallel to e1 by a unitary transformation.1 Rotations into the x1 Axis An important preliminary fact concerns vector rotations. take any orthogonal matrix S whose first column is x s1 = .5. and since all the other sj are orthogonal to s1 .3.  =  . Numerically stable and efficient algorithms exist for the Schur decomposition. so some matrices cannot be diagonalized by similarity transformations. if Ax = λx . This transformation is equivalent to factoring A into the product A = ST S H . x Since sT x = xT x/ x = x . so λ is also an eigenvalue for T . we will not study these algorithms. then the eigenvalues of the triangular matrix T are equal to those of the original matrix A.3. It is intuitively obvious that any nonzero real vector x can be rotated into a vector parallel to e1 . but only show that all square matrices admit a Schur decomposition. But the trick is the same: let s1 = x .  . as we now show. In fact.  x =  . then that is. sT sT x n n 0 which is parallel to e1 as desired. T = Q−1 AQ . . since it triangularizes any square matrix A by a unitary (possibly complex) transformation: T = S H AS . COMPUTING EIGENVALUES AND EIGENVECTORS NUMERICALLY 61 5. The eigenvectors.3 Computing Eigenvalues and Eigenvectors Numerically The examples above have shown that not every n × n matrix admits n independent eigenvectors. It is important to notice that if a matrix A is triangularized by similarity transformations. Fortunately. however. The Schur decomposition does even better. QT Q−1 x = λx .   . are changed according to the last equation.   0  ST x =  . Formally. In this note.  . x . and this product is called the Schur decomposition of A. We will show later on that this allows solving systems of linear differential equations regardless of the structure of the system’s matrix of coefficients. Let e1 be the first column of the identity matrix.

Consequently. .   SH x =  . its eigenvalues can be determined. . if the Schur decomposition of a matrix can be computed. . By substituting this into (5.   . and unitary transformations preserve eigenvalues. . The diagonal elements of a triangular matrix are its eigenvalues.  x =  .   .12) and rearranging we have     r r  0   0      AU  .   . Let U be a unitary transformation that transforms the (possibly complex) eigenvector x of A into a real vector on the x1 axis:   r  0    x=U . EIGENVALUES AND EIGENVECTORS Now s1 may be complex. 0 1 0 . . We have sH x = xH x/ x = x .  .12) . .  just about like before. 0     =      =     λ     λ  r 0 . as we will see later.3. .62 CHAPTER 5. a system of linear differential equations can be solved regardless of the structure of the matrix of its coefficients. .  T = U H AU (5. 0 Proof. 0 where r is the nonzero norm of x. 5.  = λU  .   . Lemma 5. . .  =  . Ax = λx then there is a transformation where U is a unitary. Moreover. 0 1 0 .1 If A is an n × n matrix and λ and x are an eigenvalue of A and its corresponding eigenvector. 0 0              C  . n × n matrix. such that    T =  λ 0 . .    U H AU      U H AU   0 r 0 .2 The Schur Decomposition The Schur decomposition theorem is the cornerstone of eigenvalue computations.  . .3. 0     . and 1  sH 1   sH x 1     . It states that any square matrix can be triangularized by unitary transformations. . . We are now ready to triangularize an arbitrary square matrix A.  sH sH x n n x 0 .

5. and S H AS is upper-triangular. 0 Clearly. Because we can pick any eigenvalue as λ. This is the preferred way to compute eigenvalues numerically. Fortunately. COMPUTING EIGENVALUES AND EIGENVECTORS NUMERICALLY    T  1 0 . .  63      =    The last left-hand side is the first column of T . . Then from the lemma there exists a unitary U such that   λ  0    U H AU =  . S can be chosen so that the eigenvalues λi of A appear in any order along the diagonal of T . V  . 0     . 0   λ 0 .2 (Schur) If A is any n × n matrix then there exists a unitary n × n matrix S such that S H AS = T where T is triangular.  . Proof. Since the elements on the diagonal of a triangular matrix are the eigenvalues.  . S is a unitary matrix. By induction. . and the corresponding right-hand side is of the form required by the lemma. . The theorem obviously holds for n = 1: 1A1 = A . C   .3. there is a stable and efficient algorithm to compute the Schur decomposition. Partition C into a row vector and an (n − 1) × (n − 1) matrix G: C= wH G . 0 where λ is any eigenvalue of A.  . ∆ Theorem 5. . Furthermore. there is a unitary matrix V such that V H GV is a Schur decomposition of G. By the inductive hypothesis. Let   0 ··· 0 1  0    S=U . S H AS is the Schur decomposition of A. the order of eigenvalues can be chosen arbitrarily.3. ∆ This theorem does not say how to compute the Schur decomposition. only that it exists. Suppose it holds for all matrices of order n − 1.

But eigenvectors are the solution of the homogeneous system (5. Furthermore. We already know that Hermitian matrices (and therefore real and symmetric ones) have real eigenvalues (theorem 5. Theorem 5. the proof is complete. In the process. a Hermitian matrix. so is T . Positive definite or semidefinite matrices arise in the solution of overconstrained linear systems. Since A is Hermitian. S real . but is otherwise unnecessary. The proof is very simple.1 The determinant of a matrix is equal to the product of its eigenvalues. and T H = (S H AS)H = S H AH S = S H AS = T . ∆ ⇒ A = SΛS H ⇒ A = SΛS T . xT Qx = 1 . because AT A is positive semidefinite for every A (lemma 5. so Λ is real. and this is the first equation of the theorem (the diagonal of a Hermitian matrix must be real). which is both real and rank-deficient.5). Theorem 5. First. ∆ We saw that an n × n Hermitian matrix with n distinct eigenvalues admits n orthonormal eigenvectors (corollary 5. we know from theorem 5. Notice that a positive definite matrix is also positive semidefinite.1. we know that the eigenvalues of a matrix A are equal to those of the triangular matrix in the Schur decomposition of A. because 0∗ = 0. In fact. Let now A be real and symmetric. Let now A = ST S H be the Schur decomposition of A. While this link is very important. and therefore admits nontrivial real solutions. In fact. Proof.4.6).2 (Spectral theorem) Every Hermitian matrix can be diagonalized by a unitary matrix. All that is left to prove is that then its eigenvectors are real.6 that the determinant of a triangular matrix is the product of the elements on its diagonal. S is real. we also establish a link between singular values/vectors and eigenvalues/vectors. T = S H AS. it is useful to remember that eigenvalues/vectors and singular values/vectors are conceptually and factually very distinct entities (recall figure 5. with distinct eigenvalues or not. EIGENVALUES AND EIGENVECTORS 5.4. Λ is real and diagonal. real or not. We recall that a real matrix A such that for every nonzero x we have xT Ax > 0 is said to be positive definite. so are its eigenvectors. If in addition the matrix is real. In fact.64 CHAPTER 5. It is positive semidefinite if for every nonzero x we have xT Ax ≥ 0.4 Eigenvalues/Vectors and Singular Values/Vectors In this section we prove a few additional important properties of eigenvalues and eigenvectors. Proof. now that we have the Schur decomposition. The assumption of distinct eigenvalues made the proof simple. If we recall that a unitary matrix has determinant 1 or -1. and S H = S T . we can state the following stronger result. A = AT In either case.5). has real eigenvalues and n orthonormal eigenvectors.1. that the determinants of S and S H are the same. and every real symmetric matrix can be diagonalized by an orthogonal matrix: A = AH A real. and that the determinant of a product of matrices is equal to the product of the determinants. a general relation between determinant and eigenvalues. the Schur decomposition of a Hermitian matrix is in fact a diagonalization. Thus. But the only way that T can be both triangular and Hermitian is for it to be diagonal.4. In other words.1).1. given the Schur decomposition. Thus.2). They also occur in geometry through the equation of an ellipsoid.

Recall that the eigenvalues in the Schur decomposition can be arranged in any desired order along the diagonal. The eigenvectors of A are also its singular vectors. + cn sn with But then xT Ax = xT A(c1 s1 + . Theorem 5.4. we show that if all eigenvalues λ of a real and symmetric matrix A are positive (nonnegative) then A is positive definite (semidefinite).3 establishes one connection between eigenvalues/vectors and singular values/vectors: for symmetric.4. but for arbitrary matrices. symmetric matrix is positive semidefinite iff all its eigenvalues are nonnegative. EIGENVALUES/VECTORS AND SINGULAR VALUES/VECTORS 65 in which Q is positive definite. + cn Asn ) = xT (c1 λ1 s1 + . where both Λ and S are real. positive definite matrices are associated to quadratic forms xT Qx that represent energies or second-order moments of mass or force distributions. both left and right.4. . then its eigenvalues are nonnegative.3 implies one of the two directions: If A is real. Theorem 5. A = SΛS T . If A is positive semidefinite. . .4 A real. . . To this end. The following result relates eigenvalues/vectors with singular values/vectors for positive semidefinite matrices. ∆ Theorem 5. This result can be used to introduce a less direct link. . + cn sn ) = xT (c1 As1 + . the concepts coincide. . Lemma 5. . . Since xT Ax > 0 (or ≥ 0) for every nonzero x. .2).4. symmetric.4. sn as an orthonormal basis for Rn . Proof. or at least positive semidefinite (for instance. From the previous theorem. and positive definite. Their physical meaning makes them positive definite. then xT Ax ≥ 0 for any nonzero x. so that λ ≥ 0. i But A = SΛS T with nonnegative diagonal entries in Λ is the singular value decomposition A = U ΣV T of A with Σ = Λ and U = V = S. . Since real and symmetric matrices have n orthonormal eigenvectors (theorem 5. In fact. . let x be any nonzero vector. from Asi = λsi we obtain sT Asi = sT λsi = λsT si = λ si i i i 2 =λ. symmetric. and positive semidefinite. + cn λn xT sn = λ1 c2 + .4. . . then its eigenvalues are positive. In physics. positive semidefinite matrix A are equal to its singular values. A is positive definite (semidefinite).4. Furthermore. ∆ ci = xT si . Theorem 5. we can use these eigenvectors s1 .5 AT A is positive semidefinite. . energies cannot be negative). .5. + cn λn sn ) = c1 λ1 xT s1 + . the entries in Λ are nonnegative. . + λn c2 > 0 (or ≥ 0) 1 n because the λi are positive (nonnegative) and not all ci can be zero. and in particular sT Asi ≥ 0. positive definite matrices. If the proof of that theorem is repeated with the strict inequality. Conversely. symmetric.3 The eigenvalues of a real. It is positive definite iff all its eigenvalues are positive. Proof. and write x = c1 s1 + . we also obtain that if A is real.

Bei is the i-th column of B. that is.13) implies that the i-th column of B has the same norm as the i-th row of B. from the lemma. ∆ We have seen that important classes of matrices admit a full set of orthonormal eigenvectors. for m ≤ n. The theorem below characterizes the class of all matrices with this property. If m ≥ n and A = U ΣV T is the SVD of A. For any nonzero x we can write xT AT Ax = Ax CHAPTER 5. we have AAH = AH A if and only if T T H = T H T . we first need a lemma. . The first row has first entry t11 . we have AT A = V ΣU T U ΣV T = V Σ2 V T which is in the required format to be a (diagonal) Schur decomposition with S = V and T = Λ = Σ2 . a triangular matrix T for which T T H = T H T must be diagonal. Let i = 1.7 If for an n × n matrix B we have BB H = B H B. Then. the eigenvalues of AAT are the squares of the singular values of A. AAH = ST S H ST H S H = ST T H S H and AH A = ST H S H ST S H = ST H T S H . From BB H = B H B we deduce Bx 2 = xH B H Bx = xH BB H x = B H x 2 .4. Proof. . ∆ Theorem 5. if and only if A can be diagonalized by a unitary similarity transformation. The converse is also obviously true: if T is diagonal.66 Proof. . . ∆ . which is the conjugate of the i-th row of B. and reason similarly to conclude that T must be diagonal. AAH = AH A if and only if T is diagonal. and the eigenvectors of AAT are the left singular vectors of A. n. . so the only way that its norm can be |t11 | is for all other entries in the first row to be zero. the eigenvectors of AT A are the right singular vectors of A. the i-th column of the n × n identity matrix. Similarly. We now proceed through i = 2. then T T H = T H T . Let A = ST S H be the Schur decomposition of A. However. (5. . In fact.4. and B H ei is the i-th column of B H . then for every i = 1. To prove the theorem.8 An n × n matrix is normal if an only if it commutes with its Hermitian: AAH = AH A . Thus. Since conjugation does not change the norm of a vector.4. . the norm of the i-th row of B equals the norm of its i-th column. Proof. Then. the first column of T has norm |t11 |. Proof. ∆ Theorem 5. Lemma 5. n. the equality (5. . the class of all normal matrices.6 The eigenvalues of AT A with m ≥ n are the squares of the singular values of A.13) If x = ei . the norm of the i-th row of T is equal to the norm of its i-th column. Because S is invertible (even unitary). for m ≤ n. AAT = U ΣV T V ΣU T = U Σ2 U T is a Schur decomposition with S = U and T = Λ = Σ2 . EIGENVALUES AND EIGENVECTORS 2 ≥ 0. Similarly. that is. This is the definition of a normal matrix.

Proof. so theorem 5. EIGENVALUES/VECTORS AND SINGULAR VALUES/VECTORS Corollary 5. and orthogonal matrices are all normal. but so do.5. Notice that Hermitian (and therefore also real symmetric) matrices commute trivially with their Hermitians. normal matrix must be diagonal. unitary. 67 ∆ Checking that AH A = AAH is much easier than computing eigenvectors.4. unitary (and therefore also real orthogonal) matrices: UUH = UHU = I .4. We proved this in the proof of theorem 5. for instance. .4. real symmetric. Hermitian.8.8 is a very useful characterization of normal matrices.9 A triangular. Thus.4.

68 CHAPTER 5. EIGENVALUES AND EIGENVECTORS .

Finally. . This result is based on the Schur decomposition introduced in chapter 5. in which x0 is a known vector.2). defines the initial value of the solution.2) x(0) = where x = x(t) is an n-dimensional vector function of time t. In section 6.3.  . . we show that scalar differential equations of order greater than one can be reduced to systems of first-order differential equations. . xn n−1 d y dtn−1 With this definition.   dt  . = x= .1) by introducing the n-dimensional vector     y x1  dy   .5.3) In fact.1 Scalar Differential Equations of Order Higher than One The first-order system (6. First. we set up and solve a particular differential system as an illustrative example. .Chapter 6 Ordinary Differential Systems In this chapter we use the theory developed in chapter 5 in order to solve systems of first-order linear differential equations with constant coefficients. .   .1) (6.2.1) subsumes also the case of a scalar differential equation of order n. in sections 6. These systems have the following form: ˙ x = Ax + b(t) x0 (6. Then. . such an equation can be reduced to a first-order system of the form (6. n − 1 . the dot denotes differentiation. which is numerically preferable to the more commonly used Jordan canonical form. the coefficients aij in the n × n matrix A are constant. in section 6. dt 69 for i = 0. . + c1 + c0 y = b(t) . we make this result more specific by showing that the solution to a homogeneous system is a linear combination of exponentials multiplied by polynomials in t. 6. n dt dt dt (6. we have di y dti dn y dtn = = xi+1 dxn . we recall a general result for the solution of first-order differential systems from the elementary theory of differential equations.4 and 6. possibly greater than 1. dn y dn−1 y dy + cn−1 n−1 + . . and the vector function b(t) is a function of time. The equation (6.

. . the exponential is introduced by other means. . ORDINARY DIFFERENTIAL SYSTEMS dxi (6. .        ··· 1 · · · −cn−1      .2) is given by x(t) = xh (t) + xp (t) where xh (t) is the solution of the homogeneous system ˙ x = x(0) = and xp (t) is a particular solution of ˙ x = x(0) = Ax + b(t) 0. .4) dt for i = 1. .   is the so-called companion matrix of (6. . For the scalar exponential eλt we can write a Taylor series expansion eλt = 1 + λt λ2 t2 + + ··· = 1! 2! ∞ j=0 λj t j . . In some treatments. . 0 b(t) 6. . introduced in the following. . .4). the exponential eZ of a matrix Z ∈ Rn×n is instead defined by the infinite series expansion Z Z2 e =I+ + + ··· = 1! 2! Z 1 Not ∞ j=0 Zj . 0 −c1     b(t) =    0 1 . 0 −c2 ··· ··· . we obtain the first-order system ˙ x = Ax + b(t) xi+1 = where     A=   0 0 .2 General Solution of a Linear Differential System We know from the general theory of differential equations that a general solution of system (6.3) together with the n − 1 differential equations (6. j! always. n − 1. . in calculus classes. For matrices.3) and 0 0 . 0 0 . . j! Usually1 . the exponential is defined through its Taylor series. If we write the original system (6. and the Taylor series expansion above is proven as a property. . Ax x0 The two solution components xh and xp can be written by means of the matrix exponential. . 0 −c0 1 0 .70 and x satisfies the additional n − 1 equations CHAPTER 6.1) with initial condition (6. .

This is easily verified. and xh (t) = eAt x(0) (6.5) deAt dt Thus. j! (6.5) gives deAt dt A2 t A3 t2 + + ··· 1! 2! At A2 t2 = A I+ + + ··· 1! 2! = A+ = AeAt .6) is a solution to the differential system (6. It turns out that this infinite sum converges (to an n × n matrix which we write as eZ ) for every matrix Z. so xp satisfies equation (6. At A2 t2 A3 t3 + + + ··· = 1! 2! 3! ∞ j=0 Aj tj .7) (6. In summary.1). and the general term Z j /j! is simply the matrix Z raised to the jth power divided by the scalar j!.2) we obtain w = x0 . It can be shown that this solution is unique.1) with b(t) = 0 and initial values (6.6.2). GENERAL SOLUTION OF A LINEAR DIFFERENTIAL SYSTEM 71 Here I is the n × n identity matrix. since by differentiating this expression for xp we obtain ˙ xp = AeAt 0 t e−As b(s) ds + eAt e−At b(t) = Axp + b(t) . The solution to ˙ x = Ax + b(t) with initial value x(0) = x0 is x(t) = xh (t) + xp (t) where and xp (t) = 0 (6. for any vector w.8) (6. By using the initial values (6.9) (6. we also know that a particular solution to the nonhomogeneous (b(t) = 0) equation (6.11) . Substituting Z = At gives eAt = I + Differentiating both sides of (6. the function xh (t) = eAt w satisfies the homogeneous differential system ˙ xh = Axh . (6. From the elementary theory of differential equations. we have the following result.10) xh (t) = eAt x(0) t eA(t−s) b(s) ds .2.1) is given by t xp (t) = 0 eA(t−s) b(s) ds .

801-36). This result is briefly reviewed in this section for convenience. the solution obtained in section 6. However. . pointing out that no one of them is appropriate for all situations.2. since it is based on the numerically sound Schur decomposition.1 A is Not Defective In chapter 5 we saw how to find the homogeneous part xh (t) of the solution when A has a full set of n linearly independent eigenvectors. The next section shows how to do this. Once these two extreme cases (nondefective matrix or all-coincident eigenvalues) have been handled. qn with corresponding eigenvalues λ1 . we will only point out that several methods exist for computing a matrix exponential.3. (6. we briefly review this solution. Then. or is does not.3 Structure of the Solution For the homogeneous case b(t) = 0. square. In chapter 5. However.10) can be written as a linear combination.5) requires too many terms for a good approximation. pp. Let Q = q1 · · · qn .3. λn . . .3. we have seen that if (but not only if) A has n distinct eigenvalues then it has n linearly independent eigenvectors (theorem 5. we seem to have all we need.1. of scalar exponentials multiplied by polynomials. the procedure can be carried out analytically if the functions in the right-hand side vector b(t) can be integrated. a similar result can be found through the Schur decomposition introduced in chapter 5. 6. and small perturbations in the data can change the results dramatically.1). ORDINARY DIFFERENTIAL SYSTEMS Since we now have a formula for the general solution to a linear differential system. The naive solution to use the definition (6.2 for the case of n coincident eigenvalues can be applied regardless to how many linearly independent eigenvectors exist. In the general theory of linear differential systems. in section 6.13) Two cases arise: either A admits n distinct eigenvalues.3. we show how to compute the homogeneous solution xh (t) in the extreme case of an n × n matrix A with n coincident eigenvalues. this is shown via the Jordan canonical form.2 is the same as would be obtained with the method of section 6. but we will not discuss how this is done2 . 6. no. we show a general procedure in section 6. we can be much more specific about the structure of the solution (6.7). the method of section 6.3. of this subsection and of the following one are based on notes written by Scott Cohen.1. However. expm(A) is the matrix exponential of A.3 is superior in practice. . . we do not know how to compute the matrix exponential.3 for solving a homogeneous or nonhomogeneous differential system for any. and produces a result analogous to that obtained via Jordan decomposition for the homogeneous part xh (t) of the solution. as we currently assume. However. Fortunately. constant matrix A. then it has n linearly independent eigenvectors q1 .3. defective or not. A full discussion of this matter is beyond the scope of these notes. Specifically. the solution procedure we introduce in section 6. and particularly so about the solution xh (t) to the homogeneous part. and we have shown how to find xh (t) by solving an eigenvalue problem. Cleve Moler and Charles Van Loan discuss a large number of different methods.12) (6. we have seen that matrices with coincident eigenvalues can still have a full set of linearly independent eigenvectors (see for instance the identity matrix). in the paper cited above. with constant vector coefficients. Moler and Van Loan point out that the Jordan form cannot be computed reliably.3. Nineteen dubious ways to compute the exponential of a matrix (SIAM Review. the matrix exponential (6. . vol. In a fundamental paper on the subject. For a nonhomogeneous system. 20. consider the first order system of linear differential equations ˙ x = x(0) = Ax x0 . .72 CHAPTER 6. In section 6. This procedure is based on backsubstitution. As we have done for the SVD and the Schur decomposition.3 If A is not defective.1. If the matrix has a full complement of eigenvectors. . To be sure. 4. . 2 In 3 Parts Matlab.3.9) of system (6. When the matrix A is constant.

16) represents n uncoupled. . homogeneous.12) is xh (t) = QeΛt Q−1 x(0).6) in this case becomes xh (t) = eS(T t)S x(0) = SeT t S H x(0) = SeΛt+N t S H x(0). where S is Hermitian. We recall that S is a unitary matrix and T is upper triangular with the eigenvalues of A on its diagonal. (6.3.14) by Q−1 on the right. and the solution to (6. 73 (6. regardless of the number of its linearly independent eigenvectors? In any case.14) where Λ = diag(λ1 . Using the relation x = Qy. H . differential equations yi = λi yi . The solution (6.12) is xh (t) = QeΛt Q−1 x(0).3. . How can we compute the matrix exponential in the extreme case in which A has n coincident eigenvalues. Similarly. and the consequent relation y(0) = Q−1 x(0).12) is xh (t) = SeΛt S H x(0). if it has n orthonormal eigenvectors q1 . we derived that the solution to (6. eλn t ). ˙ The solution is yh (t) = eΛt y(0). Multiplying both sides of (6. λn ) is a square diagonal matrix with the eigenvalues of A on its diagonal. This follows easily from the definition of eZ and the fact that (Q(Λt)Q−1 )j = Q(Λt)j Q−1 . .2). it should be the case that −1 eQ(Λt)Q = QeΛt Q−1 . Since Aqi = λi qi . The last equation (6.12) can be rewritten as follows: ˙ x = Ax ˙ x = QΛQ−1 x −1 ˙ Q x = ΛQ−1 x ˙ y = Λy. If A is normal. Comparing with (6. qn . if A = SΛS H . where Λ is diagonal and N is strictly upper triangular. then the solution to (6. we obtain A = QΛQ−1 . system (6. then Q is replaced by the Hermitian matrix S = s1 · · · sn . A admits a Schur decomposition A = ST S H (theorem 5. .12) becomes xh (t) = SeΛt S H x(0).3.2 A Has n Coincident Eigenvalues When A = QΛQ−1 . and eS(Λt)S H = SeΛt S H . 6. . that is. where eΛt = diag(eλ1 t . we see that the solution to the homogeneous system (6. .15) Then.6. . . . . Q−1 is replaced by S H . (6.6).16) where y = Q−1 x. STRUCTURE OF THE SOLUTION This square matrix is invertible because its columns are linearly independent. we have AQ = QΛ. . Thus we can write T as T = Λ + N.

N 2 has two nonzero superdiagonals. Suppose. the general solution to the homogeneous differential system (6.. The fact that N t is strictly upper triangular makes the computation of this matrix exponential much simpler than for a general matrix Z. Λt and N t commute: Λt N t = λ It N t = λt N t = N t λt = N t λ It = N t Λt . and the exponential reduces to a matrix polynomial. we have N j = 0 for all j ≥ n (i. We already know how to compute eΛt . Then N has three nonzero superdiagonals. Therefore.e.12) with initial value (6. then eZ1 +Z2 = eZ1 eZ2 = eZ2 eZ1 . 0 0 0 0 0 0 0 0 In general. It can be shown that if two matrices Z1 and Z2 commute. so it remains to show how to compute eN t . that is if Z1 Z2 = Z2 Z1 . Λ = λI is a multiple of the identity matrix containing the coincident eigenvalues of A on its diagonal. that N is 4 × 4. when all the eigenvalues of A coincide. N 3 has one nonzero superdiagonal. in our case. and N 4 is the zero matrix:     0 ∗ ∗ ∗ 0 0 ∗ ∗  0 0 ∗ ∗   0 0 0 ∗  2    N =   0 0 0 ∗ →N = 0 0 0 0 → 0 0 0 0 0 0 0 0     0 0 0 ∗ 0 0 0 0  0 0 0 0   0 0 0 0  4    N3 =   0 0 0 0 →N = 0 0 0 0  . This turns out to be almost as easy as computing eΛt when the diagonal matrix Λ is a multiple of the identity matrix: Λ = λI that is. . for a strictly upper triangular n × n matrix. In summary. Thus. and N is strictly upper triangular. In fact. in this case.74 CHAPTER 6. we can write eΛt+N t = eΛt eN t .17) where A = S(Λ + N )S H is the Schur decomposition of A. ∞ n−1 N j tj N j tj eN t = = j! j! j=0 j=0 is simply a finite sum. for example. ORDINARY DIFFERENTIAL SYSTEMS Thus we can compute (6.13) when the n × n matrix A has n coincident eigenvalues is given by n−1 xh (t) = SeΛt j=0 N j tj H S x0 j! (6. N is nilpotent of order n).6) if we can compute eT t = eΛt+N t .

y(0) =   . In fact. .  .3. and consider the transformed system ˙ y(t) = T y(t) + c(t) where y(t) = S H x(t) and c(t) = S H b(t) .20) The triangular matrix T can always be written in the following form:  T11 · · · · · · T1k  0 T22 · · · T2k  T = .  . . . .3 The General Case We are now ready to solve the linear differential system ˙ x = x(0) = Ax + b(t) x0 (6. The remaining nonzero blocks Tij with i < j can be in turn bundled into matrices Ri = Ti.3.  c(t) =  .21) (6. The vector c(t) can be partitioned correspondingly as follows   c1 (t)  . defective or not. 0 ··· 0 Tkk where the diagonal blocks Tii for i = 1. yk (t)]T else di (t) = 0 (an nk -dimensional vector of zeros) end Tii = λi I + Ni (diagonal and strictly upper-triangular part of Tii ) yi (t) = eλi It end. ck (t) where ci has ni entries.      (6. let A = ST S H be the Schur decomposition of A.i+1 ··· Ti.20) can then be solved by backsubstitution as follows: for i = k down to 1 if i < k di (t) = Ri [yi+1 (t). . . .  y(t) =  . .k that contain everything to the right of the corresponding Tii . . with arbitrary b(t). . . yk (t)   and for the initial values  y1 (0)   . and the same can be done for  y1 (t)  . .6.19) in the general case of a constant matrix A. STRUCTURE OF THE SOLUTION 75 6. .. yk (0) The triangular system (6. . . k are of size ni ×ni (possibly 1×1) and contain all-coincident eigenvalues.18) (6. .  . j j ni −1 Ni t j=0 j! yi (0) + t 0 eλi I(t−s) j j ni −1 Ni (t−s) j=0 j! (ci (s) + di (s)) ds .

17) with S = I. As an illustration.76 CHAPTER 6. the 2 × 2 homogeneous.18) is then x(t) = Sy(t) . we obtain y(t) = eλIt In scalar form. the expression for yi (t) is a direct application of equations (6.22) for t11 = t22 and for t11 = t22 have different forms.11). the applicability of this routine depends on whether the integral in the expression for yi (t) can be computed analytically. = t11 0 1 0 t12 t22 t12 t 1 y1 y2 . When t11 = λ1 = t22 = λ2 . we have y1 (t) = y1 (0) eλt + t12 y2 (0) t eλt t12 y2 (0) λ2 t (e − eλ1 t ) y1 (t) = y1 (0)eλ1 t + λ2 − λ1 if if t11 = t22 t11 = t22 .10). y2 (t) = t22 y2 ˙ has solution We then have and y1 (t) = y1 (0)eλ1 t + 0 t y2 (t) = y2 (0) eλ2 t . we consider a very small example. (6. since we know that in this case two linearly independent eigenvectors exist (theorem 5.22) and the initial value equation y(0) = y0 . This is certainly the case when b(t) is a constant vector b. we apply the backsubstitution procedure introduced in this section.1. ORDINARY DIFFERENTIAL SYSTEMS In this procedure. Instead. In the general case. (6.22). The solution x(t) for the original system (6.22) y(0) . y1 ˙ y2 ˙ When t11 = t22 = λ. and it is easy to verify that this solution satisfies the differential system (6. While y2 (t) is the same in both cases. (6. the solutions to system (6. we could solve the system by finding the eigenvectors of T . because then the integrand is a linear combination of exponentials multiplied by polynomials in t − s. . and (6. triangular case. which can be integrated by parts. d1 (t) = t12 y2 (t) = t12 y2 (0) eλ2 t eλ1 (t−s) d1 (s) ds t 0 = = = = y1 (0)eλ1 t + t12 y2 (0) eλ1 t y1 (0)eλ1 t + t12 y2 (0) eλ1 t 0 e−λ1 s eλ2 s ds e(λ2 −λ1 )s ds t t12 y2 (0) λ1 t (λ2 −λ1 )t e (e − 1) λ2 − λ1 t12 y2 (0) λ2 t y1 (0)eλ1 t + (e − eλ1 t ) λ2 − λ1 y1 (0)eλ1 t + Exercise: verify that this solution satisfies both the differential equation (6.1). this becomes y1 (t) = y2 (t) = (y1 (0) + t12 y2 (0) t) eλt y2 (0) eλt . The second equation of the system. Thus.9).

In fact.1. however. the new equations are differential. In the following we write the differential equations that describe this system. This would seem to present numerical difficulties when t11 ≈ t22 . because the solution would suddenly switch from one form to the other as the difference between t11 and t22 changes from about zero to exactly zero or viceversa. because of Newton’s law. By Hooke’s law.4 A Concrete Example In this section we set up and solve a more concrete example of a system of differential equations. . and since accelerations are second derivatives of position. Suppose that we want to study the evolution of the system over time. into a real solution in a form appropriate to the problem at hand.1: A system of masses and springs. This. The point of this section is to show how to transform the complex formal solution of the differential system. We will then transform these into four linear differential equations of the first order. Because differentiation occurs only with respect to one variable. the two masses would assume the positions indicated by the dashed lines. and is transformed into a first-order system with four equations. The initial system has two second-order equations. eλt − eλ1 t lim = t eλt .6. the three springs exert forces that are proportional to the springs’ elongations: f1 f2 4 Recall = = c1 v1 c2 (v2 − v1 ) that the order of a differential equation is the highest degree of derivative that appears in it. Since forces are proportional to accelerations. these are ordinary differential equations. A CONCRETE EXAMPLE 77 1 v1 rest position of mass 1 1 2 v2 2 rest position of mass 2 3 Figure 6. is not a problem. computed with any of the methods described above. 6. which allows finding eigenvalues and eigenvectors analytically with a little trick. λ1 →λ λ − λ1 and the transition between the two cases is smooth.4. Consider the mechanical system in figure 6. as opposed to partial. Two linear differential equations of the second order4 result. In the absence of external forces. time. The 4 × 4 matrix of the resulting system has an interesting structure.

the initial conditions (6. For uniformity.23) becomes u3 ˙ u4 ˙ =B u1 u2 .23). we will first transform it to a system of four first-order equations.25)  u3   v1  ˙ u4 v2 ˙ so that u3 = v1 ˙ while the original system (6. in matrix form. and u4 = v2 . ORDINARY DIFFERENTIAL SYSTEMS −c3 v2 where the ci are the positive spring constants (in newtons per meter). the trick is to introduce variables to denote the first-order derivatives of v.  (6. To solve the second-order system (6.23) . We also assume that initial conditions v(0) and ˙ v(0) (6. . v(0) ˙ v(0) .24) are replaced by the (known) vector u(0) = In the next section we solve equation (6. according to Newton’s second law: m1 v1 ¨ m2 v2 ¨ or. As shown in the introduction to this chapter.26) Likewise. ˙ We can now gather these four first-order differential equations into a single system as follows: ˙ u = Au where 0  0  A=  0 0 B 1 0 0 0 0 1 0 0    .24) are given. so that second-order derivatives of v are first-order derivatives of the new variables. The accelerations of masses 1 and 2 (springs are assumed to be massless) are proportional to their accelerations. ¨ = Bv v where v= v1 v2 and B= +c − c1m1 2 c2 m2 c2 m1 c2 +c3 − m2 = −f1 + f2 = −c1 v1 + c2 (v2 − v1 ) = −(c1 + c2 )v1 + c2 v2 = −f2 + f3 = −c2 (v2 − v1 ) − c3 v2 = c2 v1 − (c2 + c3 )v2 (6. which specify positions and velocities of the two masses at time t = 0.26). we define four new variables     u1 v1  u   v  u= 2 = 2  (6.78 f3 = CHAPTER 6.

y x= . √ √ −α + γ .2 = −α ± γ . where γ= α2 − β = 1 4 c1 + c2 c2 + c3 − m1 m2 2 + c2 2 . The eigenvalues µ1 and µ2 of B are the solutions of det(B − µI) = where α= 1 2 c1 + c2 c2 + c3 + m1 m2 and β= c1 c2 + c1 c3 + c2 c3 m1 m2 c1 + c2 +µ m1 c2 + c3 c2 2 +µ − = µ2 + 2αµ + β = 0 m2 m1 m2 are positive constants that depend on the elastic properties of the springs and on the masses. Here. (6. We also have that α ≥ γ. m1 m2 The constant γ is real because the radicand is nonnegative.27) . λ2 = − −α + γ . where we recall that A= 0 B I 0 and B= +c − c1m1 2 c2 m2 c2 m1 c2 +c3 − m2 (6.28) In other words.6. so that the two solutions µ1. the zeros in A are 2 × 2 matrices of zeros.5. = λy = λz . and I is the 2 × 2 identity matrix. With the system of springs in figure 6. z we can write Ax = 0 B I 0 y z = z By so that the eigenvalue equation (6. We then obtain µ1. the eigenvalues of A are the square roots of the eigenvalues of B: if we denote the two eigenvalues of B as µ1 and µ2 . The eigenvalues of A are solutions to the equation Ax = λx . If we partition the vector x into its upper and lower halves y and z. and the four eigenvalues of A.27) can be written as the following pair of equations: z By which yields By = µy with µ = λ2 .30) .2 are real and negative.1. (6.29) λ1 = √ √ λ3 = −α − γ . then the eigenvalues of A are √ √ √ √ λ1 = µ1 λ2 = − µ1 λ3 = µ2 λ4 = − µ2 .5 Solution of the Example Not all matrices have a full set of linearly independent eigenvectors. SOLUTION OF THE EXAMPLE 79 6. λ4 = − −α − γ (6. however. we are lucky.

Finally.34) In these expressions.33). and from equation (6. second-order system (6. and the two scalar equations in (6.32) Since ±(−α ± γ) are eigenvalues of B. y denotes the two eigenvectors of B. then there are two corresponding eigenvectors for A of the form y x= . In fact. The amplitudes and phases of the sinusoids. This completes the solution of our system of differential equations.25) by noticing that v is equal to the first two components of u. The first equation reads − c2 c1 + c2 − α ± γ y1 + y2 = 0 m1 m1 and is obviously satisfied by any vector of the form y=k c1 +c2 m1 c2 m1 −α±γ where k is an arbitrary constant. . and Q is in equation (6. the values of λi are given in equations (6.23) can be obtained from equation (6.80 CHAPTER 6. However. from equation (6.31) the four eigenvectors of A are proportional to the four columns of the following matrix:   c2 c2 c2 c2  Q=  a + λ2 1 c λ1 m21 λ1 a + λ2 1 m1 a + λ2 2 c λ2 m21 λ2 a + λ2 2 a= m1 a + λ2 3 c λ3 m21 λ3 a + λ2 3 m1 a + λ2 4 c λ4 m21 λ4 a + λ2 4 m1    (6. (6. Since we just found four distinct eigenvectors. ORDINARY DIFFERENTIAL SYSTEMS come in nonreal. we can write more simply u(t) = QeΛt Q−1 u(0) where λ1  0 Λ=  0 0  0 λ2 0 0 0 0 λ3 0  0 0   . on the other hand.28) we see that if y is an eigenvector of B corresponding to eigenvalue µ = λ2 .30).33) where c1 + c2 . it may be useful to add some algebraic manipulation in order to show that the solution is indeed oscillatory. Also the eigenvectors of A can be derived from those of B. the masses’ motions can be described by the superposition of two sinusoids whose frequencies depend on the physical constants involved (masses and spring constants). complex-conjugate pairs.26) is then given by equation (6.17). For k = 0. the solution to the original.32) must be linearly dependent. we note that u(t) = QeΛt w . As we see in the following. (6. since our system of springs obviously exhibits an oscillatory behavior. This is to be expected. To simplify our manipulation. m1 The general solution to the first-order differential system (6. the determinant of this equation is zero. depend on the initial conditions. however.31) ±λy The eigenvectors of B are the solutions of (B − (−α ± γ)I)y = 0 . 0  λ4 (6.

36) follow. by using the relation k1 − k2 is purely imaginary k3 − k4 is purely imaginary . SOLUTION OF THE EXAMPLE where we defined 81 w = Q−1 u(0) . 2 and simple trigonometry we obtain v(t) = q1 A1 cos(ω1 t + φ1 ) + q2 A2 cos(ω2 t + φ2 ) where ω1 = √ α−γ = 1 (a + b) − 2 1 (a + b) + 2 c2 1 2 (a − b)2 + 4 m1 m2 1 c2 2 (a − b)2 + 4 m1 m2 ω2 = √ α+γ = . :) = where we defined q1 = Thus. :)eΛt w . second-order problem. the ki must come in complex-conjugate pairs: and (6. ∗ k1 = k2 ∗ k3 = k4 . We have v(t) = Q(1 : 2. Numerical values for the constants can be found from the initial conditions u(0) by equation (6. we have Q(1 : 2. Since the λs are imaginary but v(t) is real. Since λ2 = −λ1 and (see equations (6. and derive the general solution v(t) for the original.30)). where Q(1 : 2. (6.35).35) We now leave the constants in w unspecified. we have v(0) = q1 (k1 + k2 ) + q2 (k3 + k4 ) and from the derivative ˙ v(t) = q1 λ1 k1 eλ1 t − k2 e−λ1 t + q2 λ3 k3 eλ3 t − k4 e−λ3 t we obtain ˙ v(0) = q1 λ1 (k1 − k2 ) + q2 λ3 (k3 − k4 ) .36) In fact.5. v(t) = q1 k1 eλ1 t + k2 e−λ1 t + q2 k3 eλ3 t + k4 e−λ3 t . Finally. this means that k1 + k2 is real k3 + k4 is real from which equations (6. ejx + e−jx = cos x . :) denotes the first two rows of Q.6. we can write c2 m1 c1 +c2 m1 + λ4 = −λ3 q1 q2 q2 c2 m1 c1 +c2 m1 + q1 λ2 1 and q2 = λ2 3 . Since the vectors qi are independent (assuming that the mass c2 is nonzero).

Re(k1 )) A2 = 2|k3 | φ2 = arctan2 (Im(k3 ). ORDINARY DIFFERENTIAL SYSTEMS b= c2 + c3 . y) = (0. CHAPTER 6. Im denote the real and imaginary part and where the two-argument function arctan2 is defined as follows for (x.25). ˙ The two constants k1 and k3 can be found from the given initial conditions v(0) and v(0) from equations (6. . on the other hand. 0).82 and a= c1 + c2 m1 . and not on the initial conditions. and extends the function by continuity along the y axis. This function returns the arctangent of y/x (notice the order of the arguments) in the proper quadrant. depend on the constants ki as follows: A1 = 2|k1 | . 0)  y if x > 0   arctan( x ) y  π + arctan( x ) if x < 0 arctan2 (y. y) = (0. m2 Notice that these two frequencies depend only on the configuration of the system. φ1 = arctan2 (Im(k1 ). Re(k3 )) where Re. The amplitudes Ai and phases φi . x) = if x = 0 and y > 0  π  2π  −2 if x = 0 and y < 0 and is undefined for (x.35) and (6.

The output can be the temperature of the reaction product. To this purpose.7. in the particular form of Kalman filtering. and in section 7. we will develop the theory of the Kalman filter. • A mass suspended from a spring. Kalman filtering has intimate connections with the theory of algebraic linear systems we have developed in chapters 2 and 3. the temperature of the gas being supplied. Simple examples of a dynamic system are the following: • An electric circuit. the term system refers to some physical entity on which some action is performed by means of an input u. whose inputs are the external temperature. One often says that a system evolves over time. as well as in most other sciences. is to determine a good model for the phenomena and events that are involved. that of modeling the trajectory of an enemy mortar shell. • A chemical reactor.5. models take on frequently the form of dynamic systems. A dynamic system is a mathematical description of a quantity that evolves over time. studying manipulation requires defining models for how a robot arm can move and for how it interacts with the world. 83 . The input is the force applied to the mass and the output is the position of the mass. the theory of state estimation. The definition is then illustrated by setting up the dynamic system equations for a simple but realistic application. When motion is involved. The system reacts to this input and produces an output y (see figure 7. In sections 7. and the supply rate of the gas.1). whose input is the current in a given branch and whose output is a voltage across a pair of nodes.Chapter 7 Stochastic State Estimation Perhaps the most important part of studying a problem in robotics or vision. Although in this chapter we will barely scratch its surface. The theory of dynamic systems is both rich and fascinating. For instance. As discussed in section 7. A dynamic system is a system whose phenomena occur over time. as is very often the case. an informal definition of a dynamic system is given in the next section.1: A general system.1 Dynamic Systems In its most general meaning.6 we will see that the shell can be shot down before it hits us. we will consider one of its most popular and useful aspects. Analyzing image motion implies defining models for how points move in space and how this motion projects onto the image. u input S system y output Figure 7.3 through 7. 7.

one can show that there exists a function f such that the state x(t) of the system satisfies the differential equation ˙ x(t) = f (x(t). what is input and what is output is a choice that depends on the application. t1 . for which the functional φ is continuous. Also. the system is said to be continuous. since the evolution is exactly determined once the initial state x at time 0 is known. In that case. the state could be the position and velocity of the mass. can be given in the case of a regular system. Similarly. t2 ). uk . called the state. Thus. so one cannot compute xk+1 as a function of xk . and the output is a function of the state. the output y of the system happens to coincide with one of the two state variables.1. u(t). the modeling problem is to somehow correlate inputs (causes) with outputs (effects). that is. For continuous systems. k) where f is some function that represents the change. 7. is essentially the same as for the discrete case: y(t) = h(x(t). k) . In general. Specifically. the laws that govern a given physical . This is. t2 ) where u(·) represents the entire function u. must allow computing the state (and hence the output) at time t2 . t) .84 CHAPTER 7. the system is said to be discrete. rather than functionals. the way that the input changes the state at time instant number k into the new state at time instant k + 1 can be represented by a simple equation: xk+1 = f (xk . For the mass attached to a spring. In fact. all quantities are vectors. an unrealistic state of affairs. all the quantities in the examples vary continuously with time. and k.t2 ) . if time varies discretely. which summarizes information about the past and the present of the system. for instance.2 Uncertainty The systems defined in the previous section are called deterministic.1. knowledge of both x(t1 ) and u[t1 . Rather. 7. as for instance for switching networks and computers. but rather compute x(t2 ) as a functional φ of x(t1 ) and the entire input u over the interval [t1 . the value x(t) taken by the state at time t must be sufficient to determine the output at the same point in time. Specifying the initial state x0 completes the definition of a continuous dynamic system. the relation between state and output can be expressed by another function: yk = h(xk . time does not come in quanta. The relation from state to output. on the other hand. In other cases. An effort of abstraction is therefore required. If time varies continuously. which leads to postulating a new quantity. A description of the system in terms of functions. uk . continuous or discrete. Furthermore. For a discrete system. in a dynamic system the input affects the state. Determinism implies that both the evolution function f and the output function h are known exactly. STOCHASTIC STATE ESTIMATION In all these examples. and uk is the input at time k. and with continuous first derivative. t2 ): x(t2 ) = φ(x(t1 ). In practice. u(·). in this example. it is more natural to consider time as a discrete variable. of the state at time t1 and the input over the interval t1 ≤ t < t2 .1 State Given a dynamic system. differentiable. however. t) where the dot denotes differentiation with respect to time. The examples above suggest that the output at time t cannot be determined in general by the value assumed by the input quantity at the same point in time. the laws of classical mechanics allow computing the new position and velocity of the mass at time t2 given its position and velocity at time t1 and the forces applied over the interval [t1 . A discrete dynamic system is completely described by these two equations and an initial state x0 . the output is the result of the entire history of the system. and is therefore always deducible from the latter. not just one of its values. Also.

that is. and the output y is in Rm . we want to determine the state x of the system. the input matrix G is n × p.1. This is a very important task. on a hill that looks about 0.5 kilometers higher than your own position. the system equations become xk+1 yk for the discrete case. In sections 7. You spotted an enemy mortar installation about thirty kilometers away. We assume that the input u is a vector in Rp . is that η and ξ are zero-mean Gaussian random variables with known covariance matrices Q and R. we consider the state estimation problem: given observations of the output y over an interval of time. and ˙ x(t) = y(t) = F (t)x(t) + G(t)u(t) + η(t) H(t)x(t) + ξ(t) = = Fk xk + Gk uk + ηk Hk xk + ξk for the continuous one. and the output matrix H is m × n. and the mean perturbations are no different. This uncertainty can be represented as noise that perturbs the equations we have presented so far. for otherwise the mean can be incorporated into the deterministic part. Then. we consider discretization issues because the physical system is itself continuous. AN EXAMPLE: THE MORTAR SHELL 85 system are known up to some uncertainty. The coefficients that appear in the equations are known only approximately. the example of the mortar shell will be discussed in order to see some of the technical issues involved in setting up the equations of a dynamic system. the state x is in Rn . You want to track incoming projectiles with a Kalman filter so you can aim your guns . and can change over time as a result of temperature changes. A more realistic model then allows for some inherent. unresolvable uncertainty in both f and h. = f (xk . t) + η(t) h(x(t). while the output is a set of observations made by a tracking system. which is sometimes valid and always simplifies the mathematics. respectively. t) + ξ(t) . the noise distributions can be assumed to have zero mean. but we choose to model it as a discrete system for easier implementation on a computer. Then.2. The mean may not be known. A discrete system then takes on the following form: xk+1 yk and for a continuous system ˙ x(t) = y(t) = f (x(t). k) + ηk = h(xk . in either f or h. u(t).2 An Example: the Mortar Shell In this section. uk . the state propagation matrix F is n × n. component wear. but this is a different story: in general the parameters that enter into the definitions of f and h must be estimated by some method. and the covariance matrix of the output noise ξ is m × m. the state is the (initially unknown) position and velocity of the shell.5. The covariance matrix Q of the system noise η is n × n. 7. In fact. 7. A common assumption.3 Linearity The mathematics becomes particularly simple when both the evolution function f and the output function h are linear. and so forth. Estimating the state then leads to enough knowledge about the shell to allow driving an antiaircraft gun to shoot the shell down in mid-flight. It is useful to specify the sizes of the matrices involved.3 through 7.7. k) + ξk Without loss of generality. In particular. the equations themselves are simple abstractions of a complex reality. For instance. in the case of the mortar shell.

5) . where I4 is the 4 × 4 identity matrix.4) You do not have very precise estimates of the noise that corrupts e and s.2) where g is the gravitational acceleration. you introduce a state update covariance matrix Q = 0. Because of this. given that you know the actual size of the projectiles used and all the camera parameters.1  ˙ 0.3) d when the projectile is at (d. you want to discretize these equations.1I4 . 7.8 × 10−3 kilometers per second squared. equal to 9. the size of the blob in pixels is s= √ 1000 . which you put along the diagonal of a 2 × 2 diagonal measurement covariance matrix R. You do not know the initial velocity of the projectiles. The projectile’s elevation is z e = 1000 (7.2) are continuous. All you have to track the shells is a camera pointed at the mortar that will rotate so as to keep the projectile at the center of the image. z). Since you do not trust your physics much. if t + dt is time instant k + 1 and t is time instant k. you are at (0. the aiming angle of the camera gives you elevation information about the projectile’s position. Similarly. you decide to ignore air drag. and dots denote derivatives with respect to time. equation (7.1) (7. and you have little time to get ready.6 kilometers/second for the horizontal component.86 CHAPTER 7.1) and (7. so you just guess some values: 0.2 seconds.6) (7. and the size of the blob tells you something about the distance. ˙ 2 The reasoning for the horizontal component d is the same.6 d  d   30     ˆ0 =  x  z  =  0. Thus. For the z component. ˙ 2 since z(0) − gt = z(t). your estimate of the initial state of the projectile is     ˙ −0. (7. d2 + z 2 (7. ˙ ˙ Consequently. where you see a blob that increases in size as the projectile gets closer. z is the vertical.1 The Dynamic System Equation Equations (7. Thus.2. 0).1 kilometers/second for the vertical component. STOCHASTIC STATE ESTIMATION accurately. except that there is no acceleration: ˙ dk+1 = dk + dk dt . you have 1 zk+1 = zk + zk dt − g(dt)2 .5 z where d is the horizontal coordinate. Since you are taking measurements every dt = 0. so you guess measurement covariances Re = Rs = 1000. you remember that the laws of motion for a ballistic trajectory are the following: d(t) z(t) ˙ = d(0) + d(0)t 1 = z(0) + z(0)t − gt2 ˙ 2 (7. From your high-school physics.2) yields z(t + dt) − z(t) 1 1 = z(0) + z(0)(t + dt) − g(t + dt)2 − z(0) + z(0)t − gt2 ˙ ˙ 2 2 1 = (z(0) − gt)dt − g(dt)2 ˙ 2 1 = z(t)dt − g(dt)2 . 0.

7) The two measurements sk and ek just defined can be collected into a single measurement vector yk = sk ek .4): sk = and after simplifying: sk = sk − ˆ 2000 dk zk ˆ ≈ −1000 dk + zk ˆ ˆ ˆ2 + z 2 (d2 + zk )3/2 ˆ2 (d2 + zk )3/2 ˆ2 d ˆ k k . and the two approximate measurement equations (7.7) and (7.8) 1000 ≈ 2 d2 + zk k 1000 ˆ d2 + zk ˆ2 k − ˆ 1000dk 1000ˆk z ˆ (d − dk ) − (z − zk ) ˆ ˆ2 + z 2 )3/2 k ˆ2 + z 2 )3/2 k (dk ˆk (dk ˆk zk zk ˆ zk ˆ ≈ 1000( − dk ) . we have zk zk − zk ˆ zk ˆ zk ˆ ˆ ek = 1000 ≈ 1000 + − (dk − dk ) . The two matrices F and G are as follows:     0 1 0 0 0  dt 1 0 0   0    F = G=  0 0 1 0   dt  . AN EXAMPLE: THE MORTAR SHELL Equations (7. ˆk ˆk ˆ dk d d d2 k so that after simplifying we can redefine the measurement to be the discrepancy from the estimated value: ek = ek − 1000 We can proceed similarly for equation (7. the control scalar u is equal to −g. To this end. From the elevation equation (7. (7.2.8) can be written in the matrix form yk = Hk xk where the measurement matrix Hk depends on the current state estimate ˆk : x   ˆ zk ˆ dk 0 0 3/2 3/2 ˆ z2 ˆ z2 (d2 +ˆk ) (d2 +ˆk )  k k Hk = −1000  zk ˆ 0 0 − d1 ˆ2 ˆ d k k (7. We want to replace these equations with linear approximations.5) and (7. the 4 × 4 matrix F depends on dt.2 The Measurement Equation The two nonlinear equations (7. we develop both equations as Taylor series around the current estimate and truncate them after the linear term. −dt2 0 0 dt 1 2 7.9) .2. ˆ ˆ ˆ dk dk d2 k (7.3) and (7.3). and the 4 × 1 control matrix G depends on dt.6) can be rewritten as a single system update equation xk+1 = F xk + Gu where  ˙ dk  d  xk =  k   zk  ˙ zk  87 is the state.7.4) express the available measurements as a function of the true values of the projectile coordinates d and z.

and R. Thus. from x0 we can simulate the system up until time t. . and perhaps predict where it will be in a few seconds. we introduce the Kalman filter from the simpler point of view of least squares estimation.” International Journal of Computer Vision. 26(4):639–656. A classical example of this situation in computer vision is the reconstruction of three-dimensional shape from a sequence of images. ηk ξk ∼ ∼ N (0. section 7. Rk ) . . July 1990. So what else is there left to do? From the observations. A single image is two-dimensional. Qk ) N (0. as they become available. section 7. ˆi|j is the estimate of the value that x x assumes at time i given the first j + 1 measurements y0 . Then. Kalman. A scheme that refines an initial estimation of the state as new observations are acquired is called a recursive1 state estimation system. Chellappa. This is a very interesting aspect of state estimation. 7. we want to know xk . and yet observations are processed one at a time.4 defines the essential notions of estimation theory that are necessary to understand the quantitative aspects of Kalman filtering. Even without noise. Most importantly. See for instance L. Kanade. one observation happens to be sufficient). Knowing the model for the mortar shell amounts to knowing the laws by which the object moves and those that relate the position of the projectile to our observations. and section 7. the second to which measurements are being used for the estimate. we do not want to have all the observations before we shoot: we would be dead by then. Also. The Kalman filter is one of the most versatile schemes for recursive state estimations. we frantically start studying state estimation.7 establishes a connection between the Kalman filter and the solution of a linear system. 82:34–45.10) (7. and T. so we can direct an antiaircraft gun to shoot down the target. a single observation yk may not be sufficient to determine the state xk (in the example. Broida. the estimation problem is defined in some more detail. The next few sections will be read under this impending threat. . In fact. the hat x means that the quantity is an estimate. . “Kalman filter-based algorithms for estimating depth from image sequences. in the hope to build a system that lets us shoot down the shell before it hits us. Matthies.5 develops the equation of the Kalman filter. at least when the dynamics of the system are known exactly (the system noise ηk is zero). The x x filter also computes an estimate Pk|k of the covariance of ˆk|k given those measurements. STOCHASTIC STATE ESTIMATION As the shell approaches us. thereby determining xk as well.” Transactions of the ASME Journal Basic Engineering. . Szeliski. and in particular Kalman filtering. . “Recursive 3-D motion estimation from a monocular image sequence. . September 1989.6 reconsiders the example of the mortar shell. T. E. . The original paper by Kalman (R. 1 The term “recursive” in the systems theory literature corresponds loosely to “incremental” or “iterative” in computer science. we would like to know where the mortar shell is right now.3 State Estimation In this section. The next section defines the state estimation problem for a discrete dynamic system in more detail. Clearly. Chandrashekhar. and R.88 CHAPTER 7. 1960) is still one of the most readable treatments of this subject from the point of view of stochastic estimation. In other words. “A new approach to linear filtering and prediction problems. Finally. 3(3):209-236. the Kalman filter computes the optimal estimate ˆk|k at time k given the measurements y0 . Section 7.J. so by itself it conveys no three-dimensional information. knowing x0 instead is equivalent. In these expressions. Kalman filters exist that recover shape information from a sequence of images.11) where the system noise ηk and the measurement noise ξk are Gaussian variables. the first k in the subscript refers to which variable is being estimated. since we have developed all the necessary tools in the first part of this course. the state of the dynamic system.” IEEE Transactions on Aerospace and Electronic Systems. It is really the ensemble of all observations that let one estimate the state. Here. . S. yj . Given a discrete dynamic system xk+1 yk = = Fk xk + Gk uk + ηk Hk xk + ξk (7. yk . ˆ as well as a (possibly completely wrong) estimate x0 of the initial state and an initial covariance matrix P0 of the estimate ˆ0 . in general.

but that we do not trust our state estimate ˆk|k−1 at all. The idea is that as time goes by the uncertainty on the state decreases.3. that is. At time k. of transforming ˆk|k−1 into ˆk|k . Now yk becomes available.3. . The update stage of the Kalman filter uses Rk and Pk|k−1 to weigh past evidence (ˆk|k−1 ) and new observations (yk ). That would mean that we trust the new measurement completely. The uncertainty about the former is Rk . just before the new measurement yk comes in. we probably need to correct our estimate of the state xk . Because a new measurement has been read.2. yk−1 . . through the measurement equation (7.7. while that about the measurements may remain the same. . we could compute the new measurement yk x x x without even looking at it. even if the latter was obtained through a large number of previous measurements. and we can consider the residue rk = yk − ˆk|k−1 = yk − Hk ˆk|k−1 . Even if ˆk|k−1 is not exact. we have an estimate ˆk|k−1 of the state vector xk based on the x previous measurements y0 . we x need some criterion for comparing the quality of the new measurement yk with that of our old estimate ˆk|k−1 of the x state. in the following sense.2: The update stage of the Kalman filter changes the estimate of the current system state xk to make the prediction of the measurement closer to the actual measurement yk . Thus. the estimate x ˆk|k−1 = Hk ˆk|k−1 y x is still our best bet. . If ˆk|k−1 were exact. This stage is represented graphically in the middle x of figure 7.1 Update The covariance matrix Pk|k must be computed in order to keep the Kalman filter running. The question however is. this uncertainty becomes usually smaller: Pk|k < Pk|k−1 .11). . y x If this residue is nonzero. by how much should we correct our estimate of the state? We do not want to make ˆk|k y coincide with yk . as well as the consequent growing uncertainty. also the uncertainty measure Pk|k−1 must be updated. measurements count less and less as the estimate approaches its true value. so that it becomes available for the next step. 7. the covariance of the observation error. STATE ESTIMATION 89 yk-1 k-1 Hk ^ xk-1 | k-1 Pk-1 | k-1 propagate ^ xk | k-1 Pk | k-1 yk k ^ yk | k-1 ^ xk | k Pk | k ^ xk+1 | k Pk+1 | k yk+1 time k+1 update propagate Figure 7. At the same time. Now we face the problem of incorporating the new measurement yk into our estimate. Propagation then accounts for the evolution of the system state. so that the new prediction ˆk|k = Hk ˆk|k y x of the measurement value is closer to the measurement yk than the old prediction ˆk|k−1 = Hk ˆk|k−1 y x that we made just before the new measurement yk was available. Then. The uncertainty about the state just before the new measurement yk becomes available is Pk|k−1 .

we see what this means.10) essentially “dead reckons” the new state from the old. just as the state vector xk represents all the information necessary to describe the evolution of a deterministic system. This increase in uncertainty depends on the system noise covariance Qk . x Hopefully. both state estimate and covariance must be propagated to the new time k + 1 to yield the new state estimate ˆk+1|k and the new covariance Pk+1|k . that is. both state estimate and state covariance matrix have been updated as described above. and again one must say in what sense ˆ is required to be “as close as possible” to x. In fact. x In summary. the Kalman filter computes the Best Linear Unbiased Estimate (BLUE) of the state. In this case. 7. (7. y = h(x) + n . . we have so far considered the criterion x that prefers a particular ˆ if it makes the Euclidean norm of the vector y − Hx as small as possible. so that the system (7. STOCHASTIC STATE ESTIMATION 7. starting with the definition of a linear estimation problem.3. In particular. and inaccuracies in our model of how this happens lead to greater uncertainty. no matter how the phrase “as close as possible” is interpreted.10). This is the x . Inverting a function is an example of estimation. The optimal estimator is then represented by the matrix L = H −1 and the optimal estimate is ˆ = Ly . The state changes according to the system equation (7. The Kalman filter itself is derived in section 7. yk should reflect this change as well. Similarly.13) is. To turn these concepts into a quantitative algorithm we need some preliminaries on optimal estimation.5. and then considering the attributes “best” and “unbiased” in turn.13) is overconstrained.90 CHAPTER 7. for a linear observation function. x A less trivial example occurs. there is usually no inverse to H. In this section.4. solving a square. in this somewhat trivial sense. then L is the inverse of h. x because of the system noise ηk . which are discussed in the next section. a problem of estimation. our uncertainty about this estimate may be somewhat greater than one time epoch ago. Both these changes are shown on the right in figure 7. nonsingular system x x y = Hx (7. For linear systems. . and any distance between ˆ and x must be zero. and its value ˆ given the observations y x x is called an estimate. the covariance matrix Pk|k contains all the necessary information about the probabilistic part of the system. in that case ˆ is equal to x. when the matrix H has more rows than columns.12) the estimation problem amounts to finding a function ˆ = L(y) x such that ˆ is as close as possible to x. 7. But between time k and time k + 1 both state and covariance may change. The function L is called an estimator. . If the function h is invertible and the noise term n is zero. and what information it needs to keep working. about how both the system noise ηk and the measurement noise ξk corrupt the quality of the state estimate ˆk|k . so our estimate ˆk+1|k of xk+1 given y0 . The system equation (7. this intuitive introduction to Kalman filtering gives you an idea of what the filter does. Thus.2. .2 Propagation Just after arrival of the measurement yk .1 Linear Estimation Given a quantity y (the observation) that is a known function of another (deterministic but unknown) quantity x (the state) plus some amount of noise.4 BLUE Estimators In what sense does the Kalman filter use covariance information to produce better estimates of the state? As we will se later.

14) are really the same problem.4.13) has no exact solution when there are more independent equations than unknowns.2. we will see that in a very precise sense ordinary least squares solve a particular type of estimation problem. This minimization problem. What the least squares approach is really saying is that even at the solution ˆ there is some residue x n = y − Hˆ x (7.7. this can be enforced by some scaling of the entries of n or. The best estimator for a linear observation function happens to be a linear estimator.13). we will probably have an estimator that produces a worse estimate than a nonlinear one.12) with h a linear function and n Gaussian zero-mean noise with the indentity matrix for covariance. the Euclidean norm of the residue (7. This amounts to saying that we have reasons to prefer the quality of some equations over others or. Notice that the observation function h can still be nonlinear.13) by a “noisy x equation”. evaluated at the solution ˆ. an overconstrained system of the form (7.2 Best In order to define what is meant by a “best” estimator. In other words.14) does not change the nature of the problem. 7. is only slightly different from its unweighted version. even. this distance is defined as the Euclidean norm of the residue vector y − Hˆ x between the left and the right-hand sides of equation (7. BLUE ESTIMATORS 91 (unweighted) least squares criterion. An estimator is said to be linear if the function L is linear. we now Wn 2 = nT R−1 n . Replacing (7. Thus. However. In fact. x 2 = nT n (the square is of course irrelevant when it comes to minimization). From the point of view of least squares.15) treats all components (all equations in (7. Even equation (7. In fact. that we want to enforce different equations to different degrees.14)) equally.14) is the correct version. if the equality sign is to be taken literally. The noise term. In the least squares approach to solving a linear system like (7. so requiring equality is hopeless. alternatively.4. called weighted least squares.15) and we would like to make that residue as small as possible in the sense of the Euclidean norm. however.13). (7.13) and its “noisy” version (7. we have W n = W (y − Hx) = W y − W Hx so we are simply solving the system W y = W Hx in the traditional. it still makes sense to look for the best possible linear estimator. nonnegative-definite matrix. y = Hx + n (7. each equation counts the same when computing the norm of the residue. However. “unweighted” sense. by some linear transformation of them: n → Wn so instead of minimizing n minimize where R−1 = W T W is a symmetric. In section 7. one needs to define a measure of goodness of an estimate. We know the solution from normal equations: ˆ = ((W H)T W H)−1 (W H)T W y = (H T R−1 H)−1 H T R−1 y .4. namely. different equations can have noise terms of different variance. can be used to generalize the problem. In fact. the estimation problem for the observation equation (7. If L is required to be linear but h is not.

Then the linear estimator L is unbiased if an only if LH = I . we give a necessary and sufficient condition for L to be unbiased. x 7. Proof. Lemma 7. which is repeated here for convenience: y = Hx + n .14). In fact. Thus.4.2. In this sense. we want to choose a linear estimator L such that the estimates ˆ = Ly x it produces minimize the following covariance matrix: P = E[(x − ˆ)(x − ˆ)T ] . and is therefore x a random vector itself.16) be zero mean. however.16) . Thus. so which norm we actually use is irrelevant. It makes therefore sense to measure the quality of an estimator by requiring that its variance be as small as possible: the fluctuations of the estimate ˆ with respect to the true (unknown) value x from one estimation experiment to the x next should be as small as possible. This x criterion is a probabilistic one. in the derivations that follow. we obtain different estimates. most interesting matrix norms are equivalent. in the sense that given two different definitions P 1 and P 2 of matrix norm there exist two positive scalars α.4.92 CHAPTER 7. requires a notion of “size” for matrices: how large is P ? Fortunately. we only use properties shared by all norms. if we estimate the same quantity many times. ˆ is a function of a random vector (noise). we can pick any norm we like. β such that α P 1 < P 2 <β P 1 . the estimates are random.1 Let n in equation (7. in the sense that if repeat the same estimation experiment many times we neither consistently overestimate nor consistently underestimate x. this translates into the following requirement: E[x − ˆ] = 0 x and E[ˆ] = E[x] . Intuitively. the identity matrix. this same solution is obtained from a completely different criterion of goodness of a solution ˆ.3 Unbiased In additionto requiring our estimator to be linear and with minimum covariance.4 The BLUE ˆ = Ly x We now address the problem of finding the Best Linear Unbiased Estimator (BLUE) of x given that y depends on x according to the model (7. We consider this different approach because it will let us show that the Kalman filter is optimal in a very useful sense. which in turn are corrupted by noise. The new criterion is the so-called minimum-covariance criterion. 7. from measurements corrupted by different noise samples from the same distribution. x x Minimizing a matrix.4. Some matrix norms were mentioned in section 3. STOCHASTIC STATE ESTIMATION Interestingly. First. Formally. The estimate ˆ of x is some function of the x measurements y. Mathematically. E[x − ˆ] x = E[x − Ly] = E[x − L(Hx + n)] = E[(I − LH)x] − E[Ln] = (I − HL)E[x] (7. we also want it to be unbiased.

2 The Best Linear Unbiased Estimator (BLUE) ˆ = Ly x for the measurement model y = Hx + n where the noise vector n has zero mean and covariance R is given by L = (H T R−1 H)−1 H T R−1 and the covariance of the estimate ˆ is x P = E[(x − ˆ)(x − ˆ)T ] = (H T R−1 H)−1 . 0 0 RLT = RR−1 H(H T R−1 H)−1 = H(H T R−1 H)−1 0 From (7. To prove that the covariance P of ˆ is given by equation (7. let L be any (other) linear unbiased estimator. that is. We can trivially write and P = = LRLT = [L0 + (L − L0 )]R[L0 + (L − L0 )]T L0 RLT + (L − L0 )RLT + L0 R(L − L0 )T + (L − L0 )R(L − L0 )T .4.18) we obtain so that (L − L0 )RLT = (L − L0 )H(H T R−1 H)−1 = (LH − L0 H)(H T R−1 H)−1 . that is. 0 But L and L0 are unbiased. = L0 RLT = (H T R−1 H)−1 H T R−1 RR−1 H(H T R−1 H)−1 0 = (H T R−1 H)−1 H T R−1 H(H T R−1 H)−1 = (H T R−1 H)−1 ∆ . the norm of the sum is greater or equal to either norm.4. when L = L0 . 0 (L − L0 )RLT = 0 . 0 the sum of two positive definite or at least semidefinite matrices.18) is the best. x x Proof.17) because L is unbiased. To show that L0 = (H T R−1 H)−1 H T R−1 L = L0 + (L − L0 ) (7. so this expression is minimized when the second term vanishes. so it is zero as well.18) is the best choice. that it has minimum covariance. This proves that the estimator given by (7. so that LH = I. Theorem 7.17). BLUE ESTIMATORS since E[Ln] = L E[n] and E[n] = 0. We can write P = = = E[(x − ˆ)(x − ˆ)T ] = E[(x − Ly)(x − Ly)T ] x x E[(x − LHx − Ln)(x − LHx − Ln)T ] = E[((I − LH)x − Ln)((I − LH)x − Ln)T ] E[LnnT LT ] = L E[nnT ] LT = LRLT (7. and The term L0 R(L − L0 )T is the transpose of this. P = L0 RLT + (L − L0 )R(L − L0 )T . 93 ∆ And now the main result.7. For such matrices. so LH = L0 H = I. For this to hold for all x we need I − LH = 0. we simply substitute L0 for L in P = LRLT : x P as promised. In conclusion.

This estimate comes with its covariance matrix Pk|k−1 . Computation occurs according to the x phases of update and propagation illustrated in figure 7.21 into one. which is related to the state xk by the equation yk = Hk xk + ξk with error covariance T E[ξk ξk ] = Rk . The filter also x computes the covariance Pk|k of the error ˆk|k − xk given those measurements. yk .5 The Kalman Filter: Derivation We now have all the components necessary to write the equations for the Kalman filter.19) where the system noise ηk and the measurement noise ξk are Gaussian random vectors. we have y = Hxk + n . One is the estimate ˆk|k−1 of the state xk given measurements up to but not x including yk . . linear. the Kalman filter computes the best. Rk ) . linear.2. (7.1 Update At time k. Thus.21) We can summarize this available information by grouping equations 7. We now apply the results from optimal estimation to the problem of updating and propagating the state estimates and their error covariances.20) The other piece of data is the new measurement yk itself. Qk ) N (0. Given a dynamic system with system and measurement equations xk+1 yk = = Fk xk + Gk uk + ηk Hk xk + ξk (7. two pieces of data are available. To summarize. unbiased estimate ˆk|k at time k given the measurements y0 . k (7.5. . . the best linear unbiased estimate ˆ of x is x where the matrix ˆ = P H T R−1 y x P = E[(ˆ − x)(ˆ − x)T ] = (H T R−1 H)−1 x x ∆ is the covariance of the estimation error. ˆ as well as the best. . 7. n ∼ N (0. unbiased estimate x0 of the initial state with an error covariance matrix P0 . ηk ξk ∼ ∼ N (0. Another way of saying this is that the estimate ˆk|k−1 differs from the true state xk by an error term ek whose covariance is Pk|k−1 : x ˆk|k−1 = xk + ek x with E[ek eT ] = Pk|k−1 . STOCHASTIC STATE ESTIMATION 7.94 CHAPTER 7.20 and 7. and packaging the error covariances into a single. R) . given a linear measurement equation y = Hx + n where n is a Gaussian random vector with zero mean and covariance matrix R. block-diagonal matrix.

because the matrices H and R contain many zeros. n= ek ξk . these two update equations are now rewritten in a more efficient and more familiar form. 95 As we know. For this reason. and the noise term ηk is zero mean. x 7.5.19. Since the new state is related to the old through the system equation 7. ˆk|k−1 x yk .2 Propagation Propagation is even simpler. THE KALMAN FILTER: DERIVATION where y= and where n has covariance R= Pk|k−1 0 0 Rk . H= I Hk . x x is usually referred to as the Kalman gain matrix. x x . We have −1 Pk|k = H T R−1 H = I T Hk −1 Pk|k−1 0 0 −1 Rk I Hk −1 T −1 = Pk|k−1 + Hk Rk Hk and ˆk|k x = = = = = In the last line. because it specifies the amount by which the residue must be multiplied (or amplified) to obtain the correction term that transforms the old estimate ˆk|k−1 of the state xk into its new x estimate ˆk|k . the difference rk = yk − Hk ˆk|k−1 x is the residue between the actual measurement yk and its best estimate based on ˆk|k−1 .7. unbiasedness requires ˆk+1|k = Fk ˆk|k + Gk uk .5. the solution to this classical estimation problem is ˆk|k x Pk|k = = Pk|k H T R−1 y (H T R−1 H)−1 . This pair of equations represents the update stage of the Kalman filter. and the matrix x T −1 Kk = Pk|k Hk Rk ∆ ∆ Pk|k H T R−1 y Pk|k −1 Pk|k−1 T −1 Hk Rk ˆk|k−1 x yk −1 T −1 Pk|k (Pk|k−1 ˆk|k−1 + Hk Rk yk ) x −1 T −1 T −1 Pk|k ((Pk|k − Hk Rk Hk )ˆk|k−1 + Hk Rk yk ) x T −1 ˆk|k−1 + Pk|k Hk Rk (yk − Hk ˆk|k−1 ) . These expressions are somewhat wasteful.

and this occurs in time for the shell to be shot down. Figure 7. 7. the dynamic system equations for a mortar shell were set up.4 shows that the distance between estimated and true target position does indeed converge to zero. we now show that estimating the initial state x0 from the first k + 1 measurements. amounts to solving a linear x system of equations with suitable weights for its rows.11) can be expanded as follows: yk = = = Hk xk + ξk = Hk (Fk−1 xk−1 + Gk−1 uk−1 + ηk−1 ) + ξk Hk Fk−1 xk−1 + Hk Gk−1 uk−1 + Hk ηk−1 + ξk Hk Fk−1 (Fk−2 xk−2 + Gk−2 uk−2 + ηk−2 ) + Hk Gk−1 uk−1 + Hk ηk−1 + ξk . obtaining ˆ0|k .7 Linear Systems and the Kalman Filter In order to connect the theory of state estimation with what we have learned so far about linear systems. ˆ0|−1 = ˆ0 x x both assumed to be given. 7. by the update equations ˆk|k x −1 Pk|k and P0|−1 = P0 .5 shows the 2-norm of the covariance matrix over time.96 CHAPTER 7.3 shows the true and estimated trajectories.3 Kalman Filter Equations ∆ ∆ In summary. Notice that coincidence of the trajectories does not imply that the state estimate is up-to-date.6 Results of the Mortar Shell Experiment In section 7. Figure 7. For this it is also necessary that any given point of the trajectory is reached by the estimate at the same time instant. x 7.2.10) and (7. = ˆk|k−1 + Kk (yk − Hk ˆk|k−1 ) x x −1 T −1 = Pk|k−1 + Hk Rk Hk where the Kalman gain is defined as and by the propagation equations T −1 Kk = Pk|k Hk Rk ˆk+1|k x Pk+1|k = = Fk ˆk|k + Gk uk x T Fk Pk|k Fk + Qk . Figure 7. The basic recurrence equations (7. STOCHASTIC STATE ESTIMATION which is the state estimate propagation equation of the Kalman filter.5. Notice that the covariance goes to zero only asymptotically. the Kalman filter evolves an initial estimate and an initial error covariance matrix. The error covariance matrix is easily propagated thanks to the linearity of the expectation operator: Pk+1|k = = = = E[(ˆk+1|k − xk+1 )(ˆk+1|k − xk+1 )T ] x x E[(Fk (ˆk|k − xk ) − ηk )(Fk (ˆk|k − xk ) − ηk )T ] x x T T Fk E[(ˆk|k − xk )(ˆk|k − xk )T ]Fk + E[ηk ηk ] x x T Fk Pk|k Fk + Qk where the system noise ηk and the previous estimation error ˆk|k − xk were assumed to be uncorrelated. Matlab routines available through the class Web page implement a Kalman filter (with naive numerics) to estimate the state of that system from simulated observations. that is.

4 0.6 1.8 0.6 0.4 1. time 2 1.3: The true and estimated trajectories get closer to one another. Trajectories start on the right.2 5 10 15 20 25 30 35 Figure 7.4 1.4: The estimate actually closes in towards the target. LINEAR SYSTEMS AND THE KALMAN FILTER 97 true (dashed) and estimated (solid) missile trajectory 1.2 1 0.7.2 1 0.2 0 −0.4 0. distance between true and estimated missile position vs.6 0.2 0 0 5 10 15 20 25 30 Figure 7.8 0.7.8 1. .

the norm of the state covariance matrix converges to zero. F1 G0 u0 + . . .22) where Φ(l. j)Gj−1 uj−1 + νk (7.23) .5: After an initial increase in uncertainty. = . Fj 1 for l ≥ j for l < j νk = Hk j=1 Φ(k − 1. . . . . F1 η0 + . Upwards segments correspond to state propagation. . We can write one system like the one in equation (7. . . . The key thing to notice about this somewhat intimidating expression is that for any k it is a linear system in x0 .   . where K is the last time instant considered. STOCHASTIC STATE ESTIMATION norm of the state covariance matrix vs time 35 30 25 20 15 10 5 0 0 5 10 15 20 25 30 Figure 7.  . downwards ones to state update. . the initial state of the system.98 40 CHAPTER 7. j) = and the term k Fl . . k yk = Hk Φ(k − 1. F0 x0 + Hk (Fk−1 . . and we obtain a large system of the form zK = ΨK x0 + gK + nK where  zK =  y0  . . K. = Hk Fk−1 Fk−2 xk−2 + Hk (Fk−1 Gk−2 uk−2 + Gk−1 uk−1 ) + Hk (Fk−1 ηk−2 + ηk−1 ) + ξk Hk Fk−1 . + Gk−1 uk−1 ) + Hk (Fk−1 . yK (7. j)ηj−1 + ξk is noise. 0)x0 + Hk j=1 Φ(k − 1. . . . + ηk−1 ) + ξk or in a more compact form.22) for every value of k = 0. .

. and weigh each equation to keep these covariances into account. + Φ(K − 1. As time goes by. . In summary. to obtain an estimate of x0 from the measurements y0 . Intuitively. The system (7.23) all at once. the Kalman filter differs from a linear solver in the following important respects: 1. having better and better approximations to the solution as new data come in is much preferable in a dynamic setting.  =  .7. . 2. K K K The least square solution to system (7. which should consequently be weighed more heavily than others.  . data my never stop arriving. HK Φ(K − 1. νK Without knowing anything about the statistics of the noise vector nK in equation (7. 1)G0 u0 + . x x 3.23) minimizes the residue between the left and the right-hand side under the assumption that all equations are to be treated the same way. 0) 0 H1 G0 u0 . The final solution can be proven to be exactly equal to solving system (7. the best we can do is to solve the system zK = ΨK x0 + gK in the sense of least squares. Better information ought to yield more accurate results. The quantitative embodiment of this intuitive idea is at the core of the Kalman filter. However.           99 nK HK (Φ(K − 1. .7. yK : ˆ0|K = Ψ† (zK − gK ) x K where Ψ† is the pseudoinverse of ΨK .23 are equally important. so that Ψ† = (ΨT ΨK )−1 ΨT . that all the noise terms in nK in equation 7. Measurements come with covariance matrices. The noise terms in nK in equation 7.23). LINEAR SYSTEMS AND THE KALMAN FILTER  ΨK   =    gK   =    H0 H1 F0 . the result with the pseudoinverse is the K same as we would obtain by solving the normal equations.23) is not solved all at once. so we ought to be able to do better. and not only for the estimate ˆ0|k of the initial state. K)GK−1 uK−1 )  ν0  . In some applications. under the assumption that the noise that affects each of the equations has the same probability distribution. and the Kalman filter makes optimal use of this information for a proper weighting of each of the scalar equations in (7. . .23). . and this is in fact the case. However. an initial solution is refined over time as new measurements become available. The Kalman filter computes up-to-date information about the current state. A solution for the estimate ˆk|k of the current state is given. we know the covariance matrices of all these noise terms. where one cannot in general wait for all the data to be collected. This is equivalent to assuming that all the noise terms in nK are equally important. . the Kalman filter for a linear system has been shown to be equivalent to a linear equation solver. . However. that is. We know that if ΨK has full rank. a small covariance means that we believe in that measurement. . knowledge of the initial state may obsolesce and become less and less useful. . and therefore in that equation. . Rather.23 are not equally important.

Sign up to vote on this title
UsefulNot useful