Algebra

I. History
The word algebra comes from the Arabic al-jabr, meaning restoration. During the 9th
century, Islamic mathematicians systemized algebra, which they called the science of
restoration and balancing. Algebraic techniques, however, have been, employed to solve
simple equations for thousands of years.

A. Ancient Origins
The ancient Egyptians and Babylonians used algebra to solve linear, quadratic, and
indeterminate equations more than 3,000 years ago. The ancient Babylonians solved
quadratic equations using essentially the same procedures taught today.
B. Babylonian Mathematical Tablet
The markings on this tablet are mathematical symbols that were, used by the Babylonian
civilization. The Babylonian system of counting was, based on the number 60 much like the
system used today is based on the number 10. The Babylonians used their system to compute
the positive roots of quadratic equations, and in geometry, they were able to find the areas
and volumes of various geometrical shapes

First century AD mathematician Hero of Alexandria and 3rd century AD mathematician

Diophantus continued the traditions of ancient Egypt and Babylon. Diophantus's book
Arithmetica is on a much higher level, however, giving many solutions to difficult
indeterminate equations. This knowledge of solutions of equations in turn found a home in
the Islamic world. In the 9th century Arab, mathematician al-Khwârizmî wrote one of the
 first systematic explanations of the basic theory of equations and included examples and
proofs. Proofs are lines of reasoning used to show that mathematical statements are true.
By the end of the 9th century, Egyptian mathematician Abu Kamil had stated and proved
the basic laws and identities of algebra and had solved complicated problems such as
finding x, y, and z such that x + y + z = 10, x2 + y2 = z2, and xz = y2.
C. Algebra in the Middle Ages and Renaissance
Ancient civilizations wrote out algebraic expressions using only occasional abbreviations.
These longhand notations were cumbersome. The exponent x6, for example, required
notation equivalent to x · x · x · x · x · x. By medieval times Islamic mathematicians were able
to talk about arbitrarily high powers of the unknown variable x, and work out the basic
algebra of polynomials (although they did not yet use modern symbolism). This included the
ability to multiply, divide, and find square roots of polynomials as well as knowledge of the
binomial theorem, which describes how to raise a binomial to an arbitrarily high power.
Persian mathematician, astronomer, and poet Omar Khayyam showed how to express roots
of cubic equations using line segments obtained by intersecting conic sections, but he could
not find a formula for the roots. A Latin translation of al-Khwārizmī's algebra text appeared in
the 12th century. In the early 13th century, the great Italian mathematician Leonardo
Fibonacci achieved a close approximation to the solution of the specific cubic equation x3 +
2x2 + cx = d. Because Fibonacci had traveled in Islamic lands, he probably used an Arabic
method of successive approximations to reach his solution.
Leonardo Fibonacci
Italian mathematician Leonardo Fibonacci made advances
in number theory and algebra. He is especially, known for his
work on series of numbers, including the Fibonacci series.
Each number in a Fibonacci series is equal to the sum of
the two numbers that came before it.Early in the 16th century,
Italian mathematicians Scipione del Ferro, Niccolò
Tartaglia, and Gerolamo Cardano found a general solution
for cubic equations. Cardano's pupil, Ludovico Ferrari, soon
found an exact solution to equations of the fourth degree,
and as a result, mathematicians for the next several centuries tried to find a formula for the
roots of equations of degree five and higher. Early in the 19th century, however, Norwegian
mathematician Niels Abel and French mathematician Évariste Galois proved that no such
formula exists.
René Descartes
The first modern philosopher, René Descartes believed science and mathematics could explain
and predict events in the physical world. Descartes developed the Cartesian coordinate system
for graphing equations and geometric shapes. Modern maps use a grid system that can be
trace back to Cartesian graphing techniques.
Hulton Deutsch
An important development in algebra in the 16th century was the introduction of modern
symbols for unknowns, algebraic powers, and algebraic operations. As a result, of this
development, Book III of La géometrie (1637), written by French philosopher and
mathematician René Descartes looks much like a modern algebra text. Descartes's most
significant contribution to mathematics, however, was his discovery of analytic geometry,
which allows problems in geometry to be, solved by solving algebraic equations. His geometry
text also contained the essentials of a course on the theory of equations, including his so-
called rule of signs for counting, the number of what Descartes called the “true” (positive) and
“false” (negative) roots of an equation. Work continued through the 18th century on the
theory of equations. In 1799, German mathematician Carl Friedrich Gauss published a proof
showing that every polynomial equation has at least one root that is a complex number.
Complex numbers are the sum of a real number and an imaginary number.

D Modern Algebra
By the beginning of the 19th century, algebra had entered its modern phase. Attention
shifted from numbers and solving polynomial equations to studying the structure of abstract
mathematical systems whose laws are, based on the behavior of mathematical objects.
Groups, sets of elements and operations that take any two elements of a set and form
another element of the set, are an example of such a system. Groups share some of the
properties of number systems but also depart from them in important ways. Groups became
one of the chief unifying concepts of 19th-century mathematics. French mathematicians
Galois and Augustin Cauchy, British mathematician Arthur Cayley, and Norwegian
mathematicians Abel and Sophus Lie made important contributions to the study of groups.
German mathematician Hermann Grassmann laid the foundations of another important
branch of modern algebra, vector analysis, during the 1840s. Vectors are mathematical
quantities that have both magnitude and direction. Despite the abstract character of vector
analysis, American physicist J. Willard Gibbs later recognized that it could be extremely useful
for physicists. The velocity of a car, for example, can be, considered a vector because it has
magnitude (speed) and direction. The widespread influence of vector analysis led British
mathematician George Boole to write An Investigation of the Laws of Thought (1854), an
algebraic treatment of basic logic. Since that time, modern algebra—also called abstract
algebra—has continued to develop. Important new results have been, discovered, and
 algebra has found applications in all branches of mathematics as well as in many of the
physical sciences.
Algebra, branch of mathematics in which symbols (usually letters) represents unknown
numbers in mathematical equations. Algebra allows the basic operations of arithmetic, such
as addition, subtraction, and multiplication, to be, performed without using specific
numbers. People use algebra constantly in everyday life, for everything from calculating how
much flour they need to bake a certain number of cookies to figuring out how long it will
take to travel by car at a certain speed to a destination that is a specific distance away.
Arithmetic alone cannot deal with mathematical relations such as the Pythagorean
Theorem, which states that the sum of the squares of the lengths of the two shorter sides of
any right triangle is equal to the square of the length of the longest side. Arithmetic can only
express specific instances of these relations. A right triangle with sides of length 3, 4, and 5, for
example, satisfies the conditions of the theorem: 32 + 42 = 52. (32 stands for 3 multiplied by itself
and is termed “three squared.”) Algebra is not limited to expressing specific instances; instead
it can make a general statement that covers all possible values that fulfill certain conditions—in
this case, the theorem: a2 + b2 = c2.
This article focuses on classical algebra, which is concerned with solving equations, uses
symbols instead of specific numbers, and uses arithmetic operations to establish ways of
handling symbols. The word algebra is also use; however, to describe various modern, more
abstract mathematical topics that also use symbols but not necessarily to represent numbers.
Mathematicians consider modern algebra a set of objects with rules for connecting or relating
them. As such, in its most general form, algebra may fairly be, described as the language of
mathematics.
Mathematical Proof is an argument that is use to show the truth of a mathematical assertion.
In modern mathematics, a proof begins with one or more statements called premises and
demonstrates using the rules of logic, that if the premises are true then a particular conclusion
must also be true.

Mathematical Proof: Figures 1 and 2

These diagrams can be used to prove the Pythagorean theorem, which states that if a right
triangle has sides of length A and B, and a hypotenuse of length C, then A2 + B2 = C2. Figure 1
and Figure 2 each have four right triangles with sides of length A and B, and a hypotenuse of
length C. Since the Figure 1 and Figure 2 both have the same area, removing the four triangles
from Figure 1 leaves a region that must have the same area as the region that is, left when the
four triangles are, removed from Figure 2. The area of the region left in Figure 1 is A2 + B2, and
the area of the region left in Figure 2 is C2. Thus A2 + B2 = C2, proving the Pythagorean Theorem
The accepted methods and strategies used to construct a convincing mathematical argument
have evolved since ancient times and continue to change. Consider the Pythagorean Theorem,
named after the 5th century BC Greek mathematician and philosopher Pythagoras, which
states that in a right-angled triangle, the square of the hypotenuse is equal to the sum of the
squares of the other two sides. Many early civilizations considered this theorem true because
it agreed with their observations in practical situations. However, the early Greeks, among
others, realized that observation and commonly held opinion do not guarantee mathematical
truth. For example, before the 5th century BC it was widely believed that all lengths could be,
expressed as the ratio of two whole numbers. However, an unknown Greek mathematician
proved that this was not true by showing that the length of the diagonal of a square with an
area of 1 is the irrational number à (see Number).An example of a mathematical proof is the
following argument, which proves that the Pythagorean Theorem is true. Figure 1 and 2
demonstrate that the relationship of A2 + B2 = C2 holds in a right-angle triangle with sides A
and B and hypotenuse C. Figure 1 shows that a square of side A + B can be divided into four of
the right-angle triangles, a square of side A, and a square of side B. Figure 2 shows that a
square of side A + B can also be dissected into four of the right-angled triangles and a square
of side C. Since the two squares of side A + B have the same area, they must still have the
same area once the four triangles are, removed from each of them. The total area of the
squares that remain on the left side is A 2 + B2, and the area of the square remaining on the
right side is C2. Thus A2 + B2 = C2.
The Greek mathematician Euclid laid down some of the conventions central to modern
mathematical proofs. His book The Elements, written about 300 BC, contains many proofs in
the fields of geometry and algebra. This book illustrates the Greek practice of writing
mathematical proofs by first clearly identifying the initial assumptions and then reasoning
from them in a logical way in order to obtain a desired conclusion. As part of such an
argument, Euclid used results that had already been, shown to be true, called theorems, or
statements that were explicitly, acknowledged to be, self-evident, called axioms; this practice
continues today.

II.
Symbol and Special terms
the symbols of algebra include numbers, letters, and signs that indicate various
arithmetic operations. Numbers are constants (values that do not change), but letters
can represent either unknown constants or variables (values that vary). Letters that are
use to represent constants are taken from the beginning of the alphabet; those used to
+ -* x = < > () { }∝≥≤±
 The basic operational signs of algebra are familiar from arithmetic: addition (+), subtraction
(-), multiplication (×), and division (÷). The multiplication symbol × is often omitted or replaced
by a dot, as in a · b. A group of consecutive symbols, such as abc indicates the product (the
result of multiplication) of a, b, and c. Division is commonly indicated by a horizontal bar (also
called a vinculum), as in: a/c. A virgule, or slash (/), may also be, used to indicate division: a/c.
A power is the product of a number multiplied by itself. The notation 42 (read, “four squared”),
for example, is used as an abbreviation for 4 · 4 (4 times 4); thus 4 2 = 16. The 4 in 42 is, called
the base, and the small raised number 2 is, called the exponent. An exponent indicates how
many times the number is multiplied by itself: x3 (read, “x cubed”) means x · x · x. More
generally xn (read “x to the nth power” or “x to the nth” where n is any number) means the
product of x multiplied by, itself n times. Fractions can take exponents as well: (1/2)2 = 1/4.
A number whose nth power is equal to x is an nth root of x. When n are 2 the term “square
root” is used and when n is 3 the term “cube root” is used. For example, 3 and -3 are both
square roots of 9 since 32 = 9 and (-3)2 = 9; 2 is a cube root of 8 since 2 3 = 8; -2 is a cube root of
-8; 1/2 is a cube root of 1/8. The square root of x is, denoted like this:
√x
The number of times the root is, multiplied by itself is, called the index. The index is usually,
omitted for square roots, but appears as a small raised number just before the root symbol for
higher roots:
√3 8=2
the two possible values of square roots, one positive and one negative are often, written using
the plus or minus symbol: ±. The equation = 2 or -2, for instance, can be abbreviated = ±2.
A. Order of operation and grouping
Algebra relies on an established sequence for performing arithmetic operations. This ensures
that everyone who executes a string of operations arrives at the same answer. Multiplication is
performed first, then division, followed by addition, then subtraction. For example:
          1 + 2 · 3
Equals 7 because 2 and 3 are multiplied first and then added to 1. Exponents and roots have
even higher priority than multiplication:
          3 · 22 = 3 · 4 = 12
Grouping symbols override the order of operations. All operations within a group are, carried
out first. Grouping symbols include parentheses ( ), brackets [ ], braces { }, and horizontal bars
that are used most often for division and roots. Adding parentheses to a previous example:

           (1 + 2) · 3

Indicates that 1 should be, added to 2 first, and then the result multiplied by 3 for a total of 9
rather than 7. Brackets and braces are use in more complicated combinations that require
multiple nested (one inside the other) groups. Operations within the innermost group are
carrying out first:
          {2[5 + 3(1 + 4)]} ={2[5 + 3 · 5]} = {2[5 + 15]} = {2 · 20} = 40
When a slash is use to indicate division, care must be taken to group the terms appropriately.
F For example,

Cannot be written ax + b/c – dy. The second notation indicates that b should be, divided by c
before b is added to ax. Grouping symbols can be used to correctly represent the fraction when
using a slash: (ax + b)/(c - dy).
B. Special definition
Any statement that contains the equality relation (=), such as 3x = 9, is called an equation. An
equation is, called an identity if the equality is true for all values of its variables; if the equation
is true for some values of its variables and false for others, the equation is conditional. The
equation x + 0 = x, for example, is an identity while 3x = 9 is conditional because it is only true
when x = 3. A term is any algebraic expression consisting only of products of constants and
variables; 2x, -a, and 1/4s4x are all examples of terms. The numerical part of a term is, called
its coefficient. The coefficients of each term above are, respectively, 2, -1, and 1/4.
An expression containing one term, such as 2x3, is, called a monomial. An expression involving
the addition or subtraction of two terms, as in 2x2 + 3x, is, called a binomial, while an
expression with three terms, such as 4x5 – x4 + 7x, is known as a trinomial. Polynomial is the
general name for expressions in which any number of terms is add or subtracted. The degree
of a polynomial refers to the largest exponent of the variables in the polynomial. For example,
if the largest exponent of a variable is 3, as in ax3 + bx2, the polynomial is said to be of degree
3. Similarly, the expression xn + xn-1 + xn-2 is of degree n.
A linear equation with one variable is a polynomial equation of degree one—that is, of the
form ax + b = 0. These are call linear equations because graphing these equations results in
straight lines. A quadratic equation in one variable is a polynomial equation of degree two—
that is, of the form ax2 + bx + c = 0.
An indeterminate equation, such as x2 + y2 = z2, involves multiple unknowns.
A prime number is any integer (the counting numbers: 1, 2, 3, …; their negatives; and zero)
that can be evenly divided only by itself and by the number 1 or the number -1. Thus, 2, 3, 5,
7, 11, 13, 17, and 19 are all prime numbers.
A factor of a number is any integer, by which the number can be, divided evenly, with no
remainder. The factors of 6, for example, are 1, 2, 3, and 6, because 6 ÷ 1 = 6, 6 ÷ 2 = 3, 6 ÷ 3 =
2, and 6 ÷ 6 = 1. The prime factors of any number are those factors to which it can be, reduced
such that the number is express only as the product of primes and their powers. For example,
the prime factors of 6 are 2 and 3. Similarly, because 60 = 22 × 3 ×5, the prime factors of 60 are
2, 3, and 5.
I. PUTTING POLYNOMIAL EQUATIONS INTO SOLVABLE FORMS
Solving a polynomial equation usually requires altering its form. The most common tools
used to manipulate equations into solvable form are the laws of arithmetic, factoring,
and the finding of least common multiples.
 A. Laws of polynomial arithmetic
In manipulating polynomials, the usual laws of the arithmetic of numbers hold. This
section lays out those laws. Normal arithmetic is limited to the set of rational numbers
(positive and negative whole numbers and infinitely repeating decimals). Algebra and
geometry can include irrational numbers (infinite decimals that do not repeat), such as pi
(symbol ∏) and the square root of 2 (abbreviated √2). The set of all rational and
irrational numbers taken together constitutes the set of real numbers.

1. Law of addition
       1. The sum of any two real numbers a and b is also, a real number, denoted a + b. The real
numbers are, closed under the operations of addition, subtraction, multiplication, division, and
the extraction of roots; this means that applying any of these operations to real numbers yields
a quantity that also is a real number.
 2. No matter how terms are, grouped in carrying out additions, the sum will always be the
same: (a + b) + c = a + (b + c). This is, called the associative law of addition.
  3. Given any real number a, there is a real number zero (0) called the additive identity, such
that a + 0 = 0 + a = a.
 4. Given any real number a, there is a number (-a), called the additive inverse of a, such that (a)
+ (-a) = 0.
 5. No matter in what order addition is, carried out, the sum will always be the same: a + b = b +
a. This is, called the commutative law of addition.
Any set of numbers obeying laws 1 through 4 is, said to form a group. If the set also obeys law
5, it is said to be an Abelian, or commutative, group-under multiplication. The set of all real
numbers, excluding zero (because division by zero is impossible), forms such a commutative
group under multiplication. Law number 4 allows the laws of multiplication to be, extended to
division, since dividing by a number is the same as multiplying by its inverse.
. Integers and real numbers are both Abelian groups. Since subtraction can be, treated as the
addition of negative numbers (3 – 4 is the same as 3 + -4), these laws also apply to subtraction.
2. Law of multiplication
Laws similar to those for addition also apply to multiplication. Since powers are a special case
of multiplication, these laws cover powers as well.
 1. The product of any two real numbers a and b is also, a real number, denoted a·b or ab.
 2. No matter how terms are grouped in carrying out multiplications, the product will always be
the same: (ab)c = a(bc). This is, called the associative law of multiplication.
 3. Given any real number a, there is a number one (1) called the multiplicative identity, such
that a(1) = 1(a) = a.
 4. Given any nonzero real number a, there is a number (a-1), or (1/a), called the multiplicative
inverse, such that a(a-1) = (a-1)a = 1.
 5. No matter in what order multiplication is, carried out, the product will always be the same:
ab = ba. This is, called the commutative law of multiplication.
Any set of elements obeying these five laws is said to be an Abelian, or commutative, group
under multiplication. The set of all real numbers, excluding zero (because division by zero is
impossible), forms such a commutative group under multiplication. Law number 4 allows the
laws of multiplication to be, extended to division, since dividing by a number is the same as
multiplying by its inverse.
3.Distributive law
Another important property of the set of real numbers links addition and multiplication in two
distributive laws as follows:
       1. a(b + c) = ab + ac
       2. (b + c)a = ba + ca
Any set of elements with an equality relation and for which two operations (such as addition
and multiplication) are, defined, and which obeys all the laws for addition, the laws for
multiplication, and the distributive laws, constitutes a field.
4.Exponent law
The laws governing exponents are extensions of the addition, multiplication, and distributive
laws. Since a2 = a · a and a3 = a · a · a, for example, a2 · a3 = a · a · a · a · a = a2+3 = a5. In general,
therefore, for real numbers:
       1. am·an = am+n
       2. (am)n = am·n
       3. am·bm = (a·b)m
B. Multiplying polynomial
To multiply polynomials, each term of each polynomial is, multiplied in turn by all of the terms
in all of the other polynomials. The results are then, added together. The following is a simple
example of the product of a binomial and a monomial: (x + 2)(4x2) = x(4x2) + 2(4x2) = 4x3 + 8x2
Each term of the one polynomial, (x + 2), is multiplied by the single term of the monomial, 4x2.
This same principle is directly, extended to polynomials of any number of terms. For example,
the product of a binomial and a trinomial is, carried out as follows:
           (5x + 2)(3x3 + x2 -4x) =
           [5x(3x3 + x2 -4x)] + [2(3x3 + x2 -4x)] =
           [5x(3x3) + 5x(x2) + 5x(-4x)] + [2(3x3) + 2(x2) + 2(-4x)] =
           [15x4 + 5x3 – 20x2] + [6x3 + 2x2 -8x] =

After such operations have been, performed, all terms of the same degree should be combined
whenever possible to simplify the entire expression:
          15x4 + 11x3 - 18x2 - 8x
 C. Factoring polynomial
Given a complicated algebraic expression, it is often useful to factor it into the product of
simpler terms. For example, 2x2 + 4xy can be factored as 2x(x + 2y) because (x + 2y) multiplied
by 2x is equal to 2x(x) + 2x(2y), which simplifies to 2x2 + 4xy. Determining the factors of a given
polynomial may be a simple matter of inspection or may require trial and error. Not all
polynomials, however, can be factored using real-number coefficients, and these are, called
prime polynomials. Some common factorizations are, given in the following examples.
 1. Trinomials of the general form x2 + (a + b)x + ab can be factored (x + a)(x + b). If a = 5 and b =
2, for example, x2 + 7x + 10 can be factored (x + 5)(x + 2). The problem can be, worked
backwards by multiplying out the factors to prove this factoring is correct:
           (x + 5)(x + 2) =
          x(x + 2) + 5(x + 2) =
          x2 + 2x + 5x + 10 =
          x2 + 7x + 10
 2. Trinomials of the general form a2x2 + (2ab)(xy) + b2y2 can be factored (ax + by)(ax + by) or
simply (ax + by)2. In the simple example where a and b are both equal to 1, x2 + 2xy + y2 can be
factored (x + y)2:
           (x + y)2 =
           (x + y)(x + y) =
          x2 + xy + xy + y2 =
          x2 +2xy + y2
Similarly, when a = 1 and b = -1, x2 – 2xy +y2 = (x – y)2. In a more complex example, where a = 5
and b = 2, 25x2 + 20xy + 4y2 can be factored (5x + 2y)2:
           (5x + 2y)2 =
           (5x + 2y) (5x + 2y) =
          25x2 + 10xy + 10xy + 4y2 =
          25x2 + 20xy + 2y2
 3. The difference of squares of the form a2x2 – b2y2 may be factored (ax + by)(ax – by). In the
simplest case, when a and b equal 1, x2 - y2 = (x + y)(x – y) because the terms of lower degree
add up to zero and thus cancel each other out:
           (x + y)(x – y) =
          x2 – xy + xy – y2 =
          x2 - y2
In the more complex case where a = 5 and b = 4, 25x2-16y2 = (5x + 4y)(5x – 4y):
           (5x + 4y)(5x – 4y) =
          25x2 – 20xy + 20xy – 16y2 =
          25x2 – 16y2
 4. The sums and differences of cubes can also be factored according to the general formulas
a3x3 + b3y3 = (ax + by)(a2x2 – axby + b2y2) and a3x3 - b3y3 = (ax - by)(a2x2 – axby + b2y2). When a and
b are both equal to 1, x3 + y3 = (x + y)(x2 - xy + y2):
           (x + y)(x2 - xy + y2) =
          x3 - x2y + xy2 + x2y - xy2 + y3 =
          x3 + y3
and x3 - y3 = (x - y)(x2 + xy + y2).
Grouping may often be useful in factoring; terms that are similar are grouped wherever
possible, as in the following example: 10x3 + 16x2y + 18xy2 -2x3 + 8x2y = 8x3 + 24x2y + 18xy2 =
2x(4x2 +12xy + 9y2) = 2x(2x+3y)2.
D. Identifying Highest Common Factor
Given a polynomial, it is frequently important to isolate the greatest common factor
from each term of the polynomial. For example, in the binomial 6x2 + 12x, the number 6
is a factor of both terms, as is x. After factoring, 6x(x + 2) is obtained, and 6x is the
greatest common factor for all terms of the original binomial. Similarly, for the trinomial
6a2x3 + 9abx + 15cx2, the number 3 is the largest numerical factor common to 6, 9, and
15, and x is the largest variable factor common to all three terms. Thus, the greatest
common factor of the trinomial is 3x and the trinomial can be factored 3x(2a2x2 + 3ab +
5cx).
E. Identifying Least Common Multiples
Find least common multiples is useful in combining algebraic fractions. The least
common multiple (LCM) of a set of numbers is the smallest number into which each
number in the set will divide evenly. The LCM of 2, 3, 4, and 6, for example, is 12.
The LCM can be, calculated by factoring numbers into their prime components. The LCM is the
product of the highest power of each prime factor of the given numbers. For example, to find
the LCM for the three numbers 27, 63, and 75, each number is first, factored: 27 = 33, 63 = 32·7,
and 75 = 3·52. The prime factors of these three numbers are 3, 5, and 7, and the highest powers
of those three factors are 33, 52, and 7. The LCM, therefore, is 33·7·52 = 4,725; 4,725 is the
smallest number into which 27, 63, and 75 will all divide evenly.
Given several algebraic expressions, the least common multiple is the expression of lowest
degree and least coefficient that can be, divided evenly by each of the expressions. To find a
common multiple of the terms 2x2y, 15x2y2, and 6ay3, all three expressions could simply be
multiplied together: (2x2y)(15x2y2)(6ay3) = 180ax4y6. However, 180ax4y6 is not the least common
multiple. To determine which is the least, each of the terms is, reduced to its prime factors. For
the numerical coefficients 2, 15, and 6, the prime factors are 2, 3·5, and 2·3, respectively; the
least common multiple for the numerical coefficients is therefore 2·3·5, or 30. Similarly,
because the constant a appears only once, it too must be a factor. Of the variables, x2 and y3
are required because they are the highest powers of the two variables that appear in any of the
expressions. The LCM of the three terms, therefore, is 30ax2y3. Each term will evenly divide this
expression:

30ax2y3 = 15y2
 2x2y

30ax2y3 = 2ay
15x2y2
30ax2y3 = 5x2
6ay3

III. Solving Equations

the means of manipulating equations outlined in the previous section can be employee to
solve equations. Given an equation, algebra supplies solutions based on the general idea of
the identity a = a. As long as the same arithmetic or algebraic procedure is, applied
simultaneously to both sides of the equation, the equality remains unaffected. The basic
strategy is to isolate the unknown term on one side of the equation and the solution on the
other.
Suppose eight people each have an equal but unknown number of pennies. The people have a
scale with 12 pennies on one side and 6 on the other. They start putting their pennies on this
unbalanced scale, trying to make it balance. The scale balances once five of the people put all
of their pennies on the side that already has six and the other three people put all of their
pennies on the side that already has 12. The number of pennies each person originally had can
be determined by turning this situation into a linear equation with one unknown:
          5x + 6 = 3x + 12
To solve this equation, the variable terms are isolated on one side and the constant terms on
the other. The term 3x can be removed from the right side by subtracting; 3x must then be
subtracted from the left side as well:
5x + 6 = 3x + 12
-3x -3x
2x + 6 12
The number 6 is then subtracted from both sides:
2x +6 = 12
-6 -6
2x 12
To isolate x on the left side, both sides of the equation are divided by 2:
2x = 6
2 2
The solution then follows directly: x = 3. Each person, then, originally had three pennies. This
can easily be verified by substituting the solution value x = 3 back into the original equation:
          5x + 6 = 3x + 12
          5(3) + 6 = 3(3) + 12
          15 + 6 = 9 + 12
          21 = 21
A. Factoring Quadratic Equation
It is not so easy to isolate the variable on one side of an equation in which more than one
degree of the variable appears. Quadratic equations are the simplest such equations. Given any
quadratic equation of the general form
          ax2 + bx + c = 0
a number of approaches are possible depending on the specific nature of the equation in
question. If the equation can be, factored, then the solution is straightforward. For instance, if
a = 1, b = -3, and c = -10, then
          x2 – 3x – 10 = 0
can be, factored as follows:
           (x – 5)(x + 2) = 0
The only way to get 0 when multiplying numbers together is when one of the numbers is 0.
Therefore this equation can only be true when one or the other of the individual factors is
equal to zero—that is, when x - 5 = 0 or x + 2 = 0. Thus, the equation has two solutions: x = 5
and x = -2. That these are the solutions to the equation may again be, verified by substituting
them back into the original equation:
          52 – 3(5) – 10 = 25 - 15 - 10 = 0
          And (-2)2 – 3(-2) - 10 = 4 + 6 - 10 = 0.
B. Completing the Square
If, on inspection, no obvious means of factoring the equation directly can be, found an
alternative might exist. For example, in the equation
          4x2 + 12x = 7
the expression 4x2 + 12x could be factored as a perfect square if it were 4x2 + 12x + 9, which
equals (2x + 3)2. This can easily be achieved by adding 9 to both sides of the equation,
completing the square:
          4x2 + 12x + 9 = 7 + 9
Factoring and simplifying yields:
           (2x + 3)2 = 16
By taking the square root of both sides, this can be, reduced to

           (2x + 3) = √16

Which is the same as?
          2x + 3 = 4 and 2x + 3 = -4
Because √16   has two values, positive 4 and negative 4. The first equation leads to the
solution x = 1/2 because 2x + 3 = 4 becomes 2x = 1 after subtracting 3 from both sides, and 2x
= 1 becomes x = 1/2 after dividing both sides by 2. The second equation leads to the solution x
= -7/2, or x = -3 1/2. Both solutions can be verified by substituting the two solutions in
question back into the original equation, 4x2 + 12x + 9 = 7 + 9:
          4(1/2)2 + 12(1/2) + 9 = 7 + 9
          4(l/4) + 6 + 9 = 16
          1 + 6 + 9 = 16
          16 = 16
          4(-31/2)2 + 12(-31/2) + 9 = 7 + 9
          4(121/4) - 42 + 9 = 16
          49 - 42 + 9 = 16
          16 = 16
C. The Quadratic Formula
Sometimes it is not possible or readily apparent how to factor an equation or complete its
square. However, all quadratic equations that can be put into the form
          ax2 + bx + c = 0
Can be solving using the quadratic formula:

−b ± √b 2−4 ac
X¿
2a

For example, to find the roots of
          x2 – 4x = -3
The equation is first, put into the standard form
          x2 – 4x + 3 = 0
In this equation a = 1, b = -4, and c = 3. These terms are then, substituted into the quadratic
formula:
−(¿−4 )± √(−22 )−4.1.3
X¿ ¿
2.1

−(−¿ 4 )± √ 4−12
X¿ ¿
2

−(¿−4 )± √ 8
X¿ ¿ ¿
2
X¿ −(−4) ± ¿)¿ 4 ± (−¿2)¿ 4+(−2) ¿ 2 ¿ 1 or
2 2 2 2
X¿ −(−4) ± ¿)¿ 4 ± (−¿2)¿ 4 – (−2 )=4 +2 ¿ 3
2 2 2 2

D. Solving Simultaneous Equations

Algebra is frequently, used to solve not just a single equation with a single unknown but also
several equations involving multiple unknowns at the same time. The problem is to find the
set of all solutions that simultaneously satisfies all of the equations. The equations to be,
solved are, called simultaneous equations, and specific algebraic techniques can be used to
solve them. For example, given the following two linear equations with two unknowns:
           (1) 3x + 4y = 10
           (2) 2x + y = 5
 When solving more than one equation at a time, we are interested in finding the set of all
solutions that will satisfy both equations. An easy way to find this set of solutions for the
linear equations 3x + 4y = 10 and 2x + y = 5 is to express one variable in terms of the other. In
this case, the variable y is isolated in the second equation, which may be rewritten y = 5 - 2x.
If we substitute this new expression of y into the first equation, we get 3x + 4(5 - 2x) = 10.
Now there is only one variable and the equation may be, solved. The solution, x = 2, may then
be substituted into both equations, which yields a value of y = 1. Thus, the set of solutions
which satisfies both equations is (2,1). Graphically, any values that satisfy both equations will
result in an intersection of a simple solution exist. The variable y in equation (2) is isolated (y =
5 - 2x), then this value of y is substituted into equation (1):
          3x + 4(5 – 2x) = 10
This reduces the problem to one involving the single linear unknown x, and it follows that
          3x + 20 – 8x = 10
Or
          -5x = -10
So 
          x = 2
When this value is, substituted into either equation (1) or (2), it follows that
          y = 1
A faster method of solving simultaneous equations is to multiply both sides of one equation by
a number that will give one of the unknown terms the same coefficient as it has in the other
equation. This multiplication allows one equation to be, subtracted from the other, completely
canceling out one of the unknowns. For example, if both sides of equation (2) are multiplied by
4, then
          8x + 4y = 20
If equation (1) is then, subtracted from equation (2) the y term cancels:
3x +4y =10
 5x = 10
The remaining equation, 5x = 10, is easy to solve, giving x = 2. This procedure leads to another
mathematical tool, matrices, which help to produce solutions for any set of linear equations
in any number of unknowns.

IV. Root (mathematics)

Root (mathematics), term used to indicate a number that when repeatedly multiplied by itself
results in a second number. For example, 2 is a third root of 8 (written, ) because 2
multiplied by itself three times yields 8 (2 × 2 × 2 = 8). Taking the root of a number is
therefore the opposite of raising a number to a power: since 3 raised to the second power is
equal to 9 (32 = 3 × 3 = 9), 3 is a second root of 9 ( = 3). Similarly, 4Â = 2 because 2 × 2 × 2 ×
2 = 16 and = 10 because (10 × 10 × 10 = 1000).
The √ symbol is, called a radical sign. The small raised number, or index, of the root is the
number of times that the root appears in the multiplication. The second root, also known as
the square root is often, written without an index: is the same as Ã. Third roots are often,
called cube roots. Roots can also be, indicated using fractions as exponents: is identical to
‚
x.
Negative numbers can be roots, and some numbers have two square roots, one positive and
one negative. For example, the cube root of –8 is –2 because –2 × –2 × –2 = 8, and the
numbers 2 and –2 are both square roots of 4 because 2 × 2 = 4 and –2 × –2 = 4.
In computation, roots obey certain rules called the laws of radicals:
(1) √n ab ¿ √n a × √n b
For example, √3 8 × √3 27 ¿ √3 8 ×27 √3 216 ¿ 6
(2) √n a ¿ ( √n a ¿
√n b b
For example, ∛ 8=∛ 8=2 /3
∛ 27 27
n m n m
(3) ( √ a ) ¿ √ a
For example,. ( √3 8 )2 ¿ √3 82 ¿ √3 64 ¿ 4
nm nk m
(4) √a ¿ √a
For example, √ 8 ¿ ¿ √2 84 ¿ √2 4096=64
3 ×2 3 4

(5) m√ √n a ¿ mn√ a
For example, √2 √3 64 ¿ 2 ×3√ 64=√6 64 ¿ 2
The root of any number with an index of 0 is undefined. The root of any number with an
index of 1 is equal to itself: .
A second mathematical use of the word “root” involves algebraic equations such as 2x = 4.
When a quantity inserted in place of the unknown quantity (x) renders the equation a true
statement, that quantity is, called a root of the equation. For 2x = 4 the root would be 2.
Logarithm, in mathematics, the exponent or power to which a stated number, called the
base, is, raised to yield a specific number. For example, in the expression 102 = 100, the
logarithm of 100 to the base 10 is 2. This is written log10100 = 2. Logarithms were originally,
invented to help simplify the arithmetical processes of multiplication, division, expansion to a
power, and extraction of a root, but they are now, used for a variety of purposes in pure and
applied mathematics.
The first tables of logarithms were publishing independently by the Scottish mathematician,
John Napier in 1614 and the Swiss mathematician Justus Byrgius in 1620. The first table of
common logarithms was compiling by the English mathematician, Henry Briggs. Common
logarithms use the number 10 as the base number. A system of logarithms often employed
uses the transcendental number as a base; they are called natural logarithms.
The method of logarithms can be illustrated by considering a sequence of powers of the
number 2: 21, 22, 23, 24, 25, and 26, corresponding to the sequence of numbers 2, 4, 8, 16, 32,
and 64. The exponents 1, 2, 3, 4, 5, and 6 are the logarithms of these numbers to the base 2.
To multiply any number in this sequence by any other number in the series it is only
necessary to add the logarithms of the numbers, then find the antilogarithm of the sum of
the logarithms, which is equal to the base number raised to the power of the sum. Thus, to
multiply 16 by 4, first note that the logarithm of 16 is 4, and the logarithm of 4 is 2. The sum
of the logarithms 4 and 2 is equal to 6, and the antilogarithm of 6 is 64, which is the product
desired. In division, the logarithms are subtracted. To divide 32 by 8 subtract 3 from 5,
giving 2, which is the logarithm of the quotient, 4.
To expand a number to any power, multiply the logarithm by the power desired and take
the antilogarithm of the product. Thus, to find 43: log 2 4 = 2; 3 × 2 = 6; antilog 6 = 64, which
is the third power of 4. Roots are, extracted by dividing the logarithm by the desired root. To
find the fifth root of 32: log232 = 5; 5 ÷ 5 = 1; antilog 1 = 2, which is the fifth root of 32.
The problem in constructing a table of logarithms is to make the intervals between
successive entries sufficiently small. In the above example, where the entries are the powers
2, 4, 8, and so on, the entries are too far apart to be useful in multiplying any larger
numbers. By advanced mathematical processes, the logarithm of any number to any base
can be, calculated, and exhaustive tables of logarithms have been prepared. Each logarithm
consists of a whole number and a decimal fraction, called respectively the characteristic and
the mantissa. In the common system of logarithms, which has the base 10, the logarithm of
the number 7 has the characteristic 0 and the mantissa .84510 (correct to five decimal
places) and is written 0.84510. The logarithm of the number 70 is 1.84510; and the
logarithm of the number 700 is 2.84510. The logarithm of the number .7 is -0.15490, which
is sometimes, written 9.84510-10 for convenience in calculation. Logarithm tables have
been, replaced by electronic calculators and computers with logarithmic functions.
Power (mathematics), term used to describe the result of repeatedly multiplying a number by
itself. The number that is successively multiplied by itself is called the base. A small raised
number called an exponent follows the base and indicates the number of times the base is to
be, multiplied. For example, 23 equals to 2 × 2 × 2, and 45 = 4 × 4 × 4 × 4× 4.
Together, a base and an exponent form an exponential expression. When an exponential
expression is, written, the base is, said to be rise to the power of the exponent. 'Ten rises to
the fourth power,' for example, would be 104 or 10,000. This phrase can, be shortened to 'ten
to the fourth power' or even 'ten to the fourth.' A number raised to the second power is
sometimes, called squared and while a number rose to the third power may be, called cubed.
The word power sometimes also means the exponent alone rather than the result of an
exponential expression.
Bases and exponents can be variables or more complicated algebraic expressions, such as (x +
y), as well as simple numbers. For example, mn means the quantity m raised to the nth power,
and (x + y)2 means (x + y)(x + y), or x2 + 2xy + y2.
Pi, Greek letter (pi) used in mathematics as the symbol for the ratio of the circumference of a
circle to its diameter. This ratio is a universal constant, meaning that it is always the same for
any size circle. Many important mathematical and physical equations contain pi. Finding the
area of a circle (A = pr2 where r is the circle’s radius) or the volume of a sphere (V = —pr3), for
instance, both require knowledge of the value of pi.
The approximate value of pi to eight decimal places is 3.14159265. The ratio is actually an
irrational number, meaning that its decimal places go on infinitely without repeating or ending
in zeros. Mathematicians have found many unending equations with solutions equal to pi; one
of the simplest is pi = 4(1 - € + ƒ - ‡ + Š …). Computers can use such equations to quickly
estimate pi with great accuracy; in this way the value of pi has been figured to more than 1.2
trillion (1,200,000,000,000) decimal places.
The ratio pi was, known in ancient times, and various approximations of its numerical value
were, used. In the Bible, the value of pi was taken to be 3. The Greek mathematician
Archimedes correctly asserted that the value was between 3 10/70 and 3 10/71. The symbol pi
for the ratio was first use in 1706 by the Welsh mathematician William Jones but it became
popular only after its adoption by the Swiss mathematician Leonhard Euler in 1737. In 1882
the German mathematician Ferdinand Lindemann proved that pi is a transcendental number—
that is, it is not the root of any polynomial equation with rational coefficients (for example, ’x3
- “x2 - 21x + 17 = 0). Consequently, Lindemann was able to demonstrate that it is impossible to
square the circle (construct a square whose area equals that of a given circle) using algebra or
a ruler and compass because the area of a square can always be expressed as a polynomial
equation with rational coefficients.
V. Parabola
Parabola is two-dimensional curve that matches the path a tossed object such as a ball
follows. Each point on the curve is equally distant from a fixed point, called the focus, and a
fixed straight line, known as the directrix. The very tip of a parabola is, called the vertex.
Parabolas are important in astronomy and physical science. An object in space follows a
parabolic orbit as it swings around a central mass if the object has just barely enough
momentum to escape from the gravity of the central mass forever. An asteroid that
followed such a path around the Sun would fly off into interstellar space, never to return.
Parabolic mirrors (reflectors that have the shape of a parabola) reflect rays of light in
parallel lines from a light source located at the mirror's focus.

Such reflectors are, used in automobile headlights and in searchlights. Parabolic mirrors
also bring parallel rays of light to a focus. This type of reflector is therefore valuable in
astronomical telescopes. Parabolic reflectors also are, used as antennas in radio astronomy
and radar.
Any parabola is symmetrical about a line that passes through the focus and is perpendicular
to the directrix, meaning that the half-parabolas on each side of the line are mirror images
of one another. A parabola can be, drawn on xy axes by graphing its equation. For a
parabola with a horizontal directrix and a vertex at (h, k), the equation is (x – h)2 = 2p(y – k),
in which p is the distance between the focus and the directrix. Conversely, the equation of a
parabola with a vertical directrix is (y – k)2 = 2p(x – h).
Parabolas are one of the conic section curves formed by the intersection of a right circular
cone and a plane. A parabola is, formed when the plane is parallel to a straight line drawn
on the slanting surface of the cone from the tip of the cone to its base.
VI. Equation
 Equation is statement of equality between two expressions, used in almost all branches of
pure and applied mathematics and in the physical, biological, and social sciences. An
equation usually involves one or more unknown quantities, called variables or in-
determinates. These are commonly denoted by letters or other symbols, as in the equations
x2 + x - 4 = 8, y = sin x + x, and 3y = log x. An equation is, named for the number of variables
it contains, called an equation in one, two, three, or more variables.
An equation is said to be satisfied or to be true for certain values of the variables if, when
the variables are replaced by these values, the expression on the left side of the equals sign
is equal to that on the right side. For example, the equation 2x + 5 = 13 is satisfied when x =
4. If one or more values of the variable fail to satisfy the equation, the equation is, called
conditional. The equation in two variables 3x + 4y = 8 is a conditional equation because it is
not satisfied when x = 1 and y = 3. An equation is, called an identity if it is, satisfied by all
possible values of the variables. For example, the equations (x + y)2 = x2 + 2xy + y2 and sin2x +
cos2x = 1 are identities because they are both true for all possible values of the unknowns. A
solution of a conditional equation is a value of the variable, or a set of values of the
variables, that satisfies the equation; thus, 3 is a solution of the equation x2 - 2x = 3; and x =
2, y = 4 is a solution of the equation 3x2 + 4y = 28. A solution of an equation in one variable is
commonly, called a root of the equation. A polynomial equation has the form a0xn + a1xn-1 +
a 2xn-2 + ... + an-2x2 + an-1x + an = 0in which the coefficients a0, a1, ..., an are constants, the
leading coefficient a0 is not equal to zero, and n is a positive integer. The greatest exponent
n is the degree of the equation. Equations of the first, second, third, fourth, and fifth
degrees are often called, respectively, linear, quadratic, cubic, biquadratic or quartic, and
quintic equations. Other important types of equations are algebraic, as in = = 7;
trigonometric, as in sin x + cos 2x = y; logarithmic, as in log x + 2 log (x + 1) = 8; and
exponential, as in 3x + 2x - 5 = 0. Diophantine equations are equations in one or more
unknowns, usually with integral coefficients, for which integral solutions are, sought.
Differential and integral equations, which involve derivatives or differentials, and integrals,
occur in calculus and its applications. A system of simultaneous equations is a set of two or
more equations in two or more unknowns. A solution of such a system is a set of values of
the unknowns that satisfies every equation of the set simultaneously.
VII. Factor (mathematics)
Factor (mathematics), any one of two or more numbers multiplied together to form a
product. The factors of any given whole number (numbers such as 1, 2, and 3 that can be
evenly, divided by 1) are themselves whole numbers by which the given number may be
divided evenly—that is, without remainder. For example, since 6 is the product of 2 × 3, 2
and 3 are factors of 6. The decimal number 1.5 is not a factor of 6, because it is not a whole
number, and 5 is not a factor of 6 because 6 cannot be divided evenly by 5. Since 1 × 6 = 6, 1
and 6 are also factors of 6. All whole numbers are factors of themselves, and 1 is a factor of
all whole numbers.
 Some numbers such as 1, 2, 3, 5, and 7, have no factors other than themselves and one;
these are prime numbers. Numbers that do have other factors are, called composite
numbers. The first five composite numbers are 4, 6, 8, 9, and 10.
The process of breaking a number down into its factors other than 1 and itself is called
factoring. For example, the composite number 72, can be factored 9 × 8. The numbers 9 and
8 are themselves composite numbers and can be further broken down into 3 × 3 and 2 × 2 ×
2. Since 3 and 2 are prime and cannot be factored further, they are called the prime factors
of 72. All composite numbers can be, expressed as products of prime numbers. Factoring
very large numbers is difficult, but the problem is of current interest because of its use in
designing secure encryption (see Cryptography).
Two or more numbers may share several factors. For example, both 18 and 24, can be
divided by 1, 2, 3, or 6. These numbers are, called the common factors of 18 and 24. The
number 6 is the largest of these common factors and is therefore, called the highest
common factor (HCF) of 18 and 24. The highest common factor is, used in a number of
important mathematical operations. For example, all fractions can be, reduced to their
lowest terms by dividing both numerator and denominator by their highest common factor:

In addition, when two or more fractions must be, multiplied, extensive calculations can
often be, avoided by finding the HCF of each fraction. For example:

In algebra, expressions can be, factored like numbers. The factors of an algebraic expression
are themselves expressions that, when multiplied together, produce the original expression.
For example, the factors of the expression x2 – y2 are (x - y) and (x + y) because (x - y)(x + y) =
x2 – y2. These are the only factors that can be, found for x2 - y2. Factoring is often an
indispensable step in solving algebraic problems.
VIII. Function
Function, in mathematics, term used to indicate the relationship or correspondence
between two or more quantities. The term function was first use in 1637 by the French
mathematician René Descartes to designate a power xn of a variable x. In 1694, the German
mathematician Gottfried Wilhelm Leibniz applied the term to various aspects of a curve,
such as its slope. The most widely used, meaning until quite recently was defined in 1829 by
the German mathematician Peter Dirichlet. Dirichlet conceived of a function as a variable y,
called the dependent variable, having its values fixed or determined in some definite
manner by the values assigned to the independent variable x, or to several independent
variables x1, x2, ..., xk.
The values of both the dependent and independent variables were real or complex
numbers. The statement y = f(x), read “ y is a function of x,” indicate the interdependence
 between the variables x and y; f(x) was usually, given as an explicit formula, such as f(x) = x2 -
3x + 5, or by a rule stated in words, such as f(x) is the first integer larger than x for all x's that
are real numbers (see Number). If a is a number, then f(a) is the value of the function for the
value x = a. Thus, in the first example, f(3) = 32 - 3 · 3 + 5 = 5, f(-4) = (-4)2 - 3(-4) + 5 = 33; in
the second example, f(3) = f(3.1) = f(p) = 4.
The emergence of set theory first extended and then altered substantially the concept of a
function. The function concept in present-day mathematics may be illustrated as follows.
Let X and Y be two sets with arbitrary elements; let the variable x represent a member of
the set X, and let the variable y represent a member of the set Y. The elements of these two
sets may or may not be numbers, and the elements of X are not necessarily of the same
type as those of Y. For example, X might be the set of the 50 states of the United States and
Y the set of positive integers. Let P be the set of all possible ordered pairs (x, y) and F a
subset of P with the property that if (x1, y1) and (x2, y2) are two elements of F, then y1≠y2
implies that x1≠x2—that is, F contains no more than one ordered pair with a given x as its
first member. (If x1≠x2, however, it may happen that y1 = y2.) A function is now regarded as
the set F of ordered pairs with the stated condition and is written F:X→Y. The set X1 of x's
that occur as first elements in the ordered pairs of F is called the domain of the function F;
the set Y1 of y's that occur as second elements in the ordered pairs is called the range of the
function F. Thus, {(New York, 7), (Ohio, 4), (Utah, 4)} is one function that has X = the set of
the 50 U.S. states and Y = the set of all positive integers; the domain is the three states
named, and the range is 4, 7.
The modern concept of a function is, related to the Dirichlet concept. Dirichlet regarded y =
x2 - 3x + 5 as a function; today, y = x2 - 3x + 5 is thought of as the rule that determines the
correspondent y for a given x of an ordered pair of the function; thus, the preceding rule
determines (3, 5), (-4, 33) as two of the infinite number of elements of the function.
Although y = f(x) is still used today, it is better to read it as “y is functionally related to x”.
A function is also called a transformation or mapping in many branches of mathematics. If
the range Y1 is a proper subset of Y (that is, at least one y is in Y but not in Y1), then F is a
function or transformation or mapping of the domain X1 into Y; if Y1 = Y,F is a function or
transformation or mapping of X1 onto Y.
IX. Geometric Progression
Geometric Progression, in mathematics, sequence of numbers in which the ratio of any term,
after the first, to the preceding term is a fixed number, called the common ratio. For
example, the sequence of numbers 2, 4, 8, 16, 32, 64, 128 is a geometric progression in
which the common ratio is 2, and 1, €, Š, }, ‰, |, ...€i, ... is a geometric progression in which
the common ratio is. The first is a finite geometric progression with seven terms; the second
is an infinite geometric progression. In general, a geometric progression may be described by
denoting the first term in the progression by a, the common ratio by r, and, in a finite
progression, the number of terms by n. A finite geometric progression may then be, written
formally as
 a, ar1, ar2 , ar3 , ar4, …, arn-1
And an infinite geometric progression as
a, ar1, ar2, …, ari, …
In general, if the nth term of a geometric progression is denoted by an, it follows from the
definition that
an =arn-1
If the symbol Sn denotes the sum of the first n terms of a geometric progression, it can be,
proved that
Sn = a(1-rn)
1-r
The terms in a geometric progression between ai, and aj, i<j, are called geometric means. The
geometric mean between two positive numbers x and y is the same as the mean
proportional Ë between the two numbers. In particular, an is the geometric mean or mean
proportional between an - 1 and an + 1.
The formal sum of the terms of an infinite geometric progression, written as
a1 + a2 + a3 + a4 + a5 + …+ an-1 + an
is called a geometric series (see Sequence and Series). In analysis, it can be proved that a
geometric series converges if the absolute value of the common ratio is less than 1;
otherwise, the series diverges. If the series does converge, the limit, S, can be shown to be
s =¿ a
1-r
The symbol is, read “the limit of Sn as n increases without bound.”Geometric series and
geometric progressions have many applications in the physical, biological, and social
sciences, as well as in investments and banking. Many problems in compound interest and
annuities are easily, solved using these concepts.
X. Graph
Graph is a diagram that shows relationships between numbers. Graphs arrange numerical
information into a picture from which it is often possible to see overall patterns or trends in
the information.
Figure 1: Point Graph
Point graphs provide a visual representation of selected data. This point graph charts the
number of glasses of lemonade sold on each day for a week. The horizontal coordinate of
each point corresponds to the number of the day. The vertical coordinate corresponds to the
number of glasses sold, as demarcated on the left hand side of the chart.
 The graph
in figure 1
illustrates
a sales
trend. It
shows the
number of glasses of lemonade sold each day for a week. To find the number of glasses sold
on day 3, first locate the number 3 on the horizontal axis and then find the point directly
above it. The position of this point corresponding to the vertical axis is 10, meaning that 10
glasses were sold on day 3. On day 1, the number of glasses sold is hard to determine
precisely, but is somewhere between 15 and 20. Though this graph, like most graphs, is not
as accurate as a list of numbers, it more clearly illustrates the overall trend that lemonade
sells better toward the end of the week (days 6 and 7) than in the middle of the week. Figure
2: Graphical Solutions
Graphs provide a valuable resource for solving simultaneous equations. Instead of solving
two related equations mathematically, it is often possible to graph the two equations to find
their point of intersection. This graph reveals that when x=2 and y=6, the equations y=3x and
y=x+4 intersect. Therefore, (2,6) is the solution to these equations
The graph in figure 2 illustrates numerical relationships. Suppose it is known that Yolanda is
four years older than Xavier. Using y for Yolanda's age and x for Xavier's age, this relationship
can be written as y = x + 4. Since 5 = 1 + 4, one possible pair of values for x and y is x = 1 and
y = 5, which can be written briefly as (1,5) (see Coordinate System). The set of all the pairs
(x,y) for which y = x + 4 is represented by the blue straight line plotted in figure 2.
Figure 3: Graph of an Inequality
The parabolic curve in this graph consists of all the points in the plane that satisfy the
equation y = x2 - 1. The shaded region inside the parabola consists of the points in which y is
greater than x2 - 1.

Graphs can be, used to solve equations simultaneously. Suppose that in addition to knowing
that Yolanda is four years older than Xavier is also know that Yolanda is three times Xavier's
age. The problem, then, is to find values for x and y that make the equations y = x + 4 and y
= 3x both true at the same time. In figure 2, these two equations are plotted together; the
solution of these simultaneous equations is the point at which the two graphs intersect
(2,6), which shows that Xavier is two years old and Yolanda is six years old.
Graphs can also be, used to exhibit inequalities. The curve in figure 3 graphs the parabola
y = x2 - 1. The shaded area, not including this curve, is the graph of the inequality y>x2 - 1.
XI. Indeterminate Equation
In mathematics, an equation that has no set of unique solutions and that therefore cannot
be uniquely, solved. An indeterminate equation may have an infinite number of solutions.
Such equations become solvable in the ordinary sense only when additional restrictions are,
put on the problem; a common restriction is that the solutions must be, expressed in
integers.
A simple example of this type of problem is the following: In what ways can change be, given
for a half-dollar in nickels and quarters? Algebraically this reduces to the equation 5x + 25y =
50. This equation has an infinite number of solutions if fractional solutions are admissible,
but the phrasing of the problem prohibits such solutions because a third of a nickel, for
example, is meaningless. With this restriction, it is apparent that three and only three
solutions exist: ten nickels and no quarters, five nickels and one quarter, or no nickels and
two quarters. Some such problems have no solutions, for example, the following one: In
what ways can change be, given for 37 cents in nickels and quarters?
In more complex, problems the solution or solutions are not apparent, and an extensive
algebra has been developed to determine the solutions. The simplest such problems can be,
expressed in the form of a linear algebraic equation in two unknowns (such as the equation
shown above) and can be, solved by a method developed by the ancient Greek
mathematicians Diophantus and Euclid. The solutions, if any exist, are determined by finding
the largest common divisor of the numbers by which x and y are multiplied in the equation;
in the equation shown above, the numbers are 5 and 25, and their largest common divisor is
5. If this largest common divisor divides the number on the right side of the equation (as 5
divides 50), the equation has one or more integral solutions.
For some further simple examples of indeterminate problems, see Diophantine Analysis.
Many of the greatest mathematicians, such as Carl Friedrich Gauss, a German, spent much
time developing integer solutions of exceedingly complex indeterminate equations. See also
Fermat's Last Theorem.

XII. Arithmetic
Arithmetic is branch of mathematics that arises from counting, the most basic mathematical
operation. Arithmetic encompasses various ways of counting, or manipulating numbers:
addition, subtraction, multiplication, and division. The ancient Greek word arithmētikē
combined the words arithmos, meaning “number,” and technē, referring to an art or skill.
Thus, arithmetic means the art of numbers. The numbers used in arithmetic may be positive
integers (whole numbers), negative integers, fractions, or decimals. For the history of
arithmetic and mathematics, see Mathematics.
Different civilizations throughout history have developed different kinds of number systems.
Although the ancient Babylonians used a system based on the number 60, all modern
cultures employ a system in which objects are, counted in groups of ten, probably because
humans have ten fingers and tend to use them in counting. This system is, called the decimal
or base 10, system.
XIII. OPERATIONS WITH POSITIVE INTEGERS
We generally count with positive integers. This unending sequence of whole numbers starts
with 1. Each subsequent number in the sequence is one more than the number before: 1, 2,
3, 4 …. Alternating numbers starting with 1 (1, 3, 5 …) are called odd numbers, while every
other number starting with 2 (2, 4, 6 …) is called an even number.
In numbers with two or more digits, every digit has a so-called place value. In the decimal
system, the place value increases from units, or ones, to tens, hundreds, thousands, and
higher as the number of digits increases from right to left. We can also say that each place
increases by a power of 10. A power of a number is the number of times the number is,
multiplied by itself. In the number 1,111, for example, the place value on the far right is a
unit or ones value; the place value just to its left is 10 (1 × 10); the next one to the left is 10 ×
10, or 100; and the place value on the far left is 10 × 10 × 10, or 1,000. The number 2,534,
then, is equivalent to (2 × 1,000) + (5 × 100) + (3 × 10) + (4 × 1).

A. Adding positive integers

The arithmetic operation of addition is basically, a means of counting quickly and is,
indicated by the plus sign (+). We could place 4 apples and 5 more apples in a row, then,
count them individually from 1 to 9. Addition, however, makes it possible to count all of the
apples in a single step (4 + 5 = 9).
We call the end, result of addition the sum. The simplest sums are usually, memorized. This
table shows the sums of any two numbers between zero and nine:

To find the sum of any two numbers from 0 to 9, locate one of the numbers in the vertical
column on the left side of the table and the other number in the horizontal row at the top.
The sum is the number, in the body of the table that lies at the intersection of the column
and row that have been, selected. For example, 6 + 7 = 13.
We can easily add long lists of numbers with more than one digit by repeatedly adding one
digit at a time. For example, if the numbers 27, 32, and 49 are, listed in a column so that all
the units are in a line, all the tens are in a line, and so on, finding their sum is relatively
simple:
First add the units (7 + 2 + 9); they total 18. Then add the digits in the tens place (2 + 3 + 4);
they total 9, but this means 9 tens, or 90. In the last step, add the total of the units to the
total of the, tens:

We can skip the second step, adding the sum of the units to the sum of the tens, by using a
shortcut called carrying. Carry the 1 in 18, which stands for 1 ten, over to the tens column
and add it directly to the digits there:

Add the digits in the tens column, including the carried 1, and place the sum, 10, just to left
of the units sum. The result is 108. Similarly, when adding numbers with three or more
places, we can carry digits to the hundreds place, thousands place, or beyond.

B. Subtracting Negative Integers

The arithmetic operation of subtraction is the opposite of addition and is, indicated by the
minus sign (-). If we take 5 apples away from 9 apples, subtraction tells how many apples
remain without our actually counting them. The simple sums memorized for addition are,
used in reverse for subtraction. For example, the result of 9 minus 5 is 4 because 4 is the
number we would have to add to 5 for a sum of 9. The end, result of subtraction is, called the
difference.
It is possible to subtract 23 from 66 by counting backward 23 integers from 66, one number
at a time, or by taking away 23 items from a collection of 66 and counting the remainder.
Either way we would reach 43. The rules of arithmetic for subtraction, however, provide a
much quicker method for obtaining the answer. We can subtract large numbers by
repeatedly subtracting one digit at a time. First align the numbers under one another, units
under units, tens under tens, as in addition:

Subtract the units: 6 - 3 = 3. Then subtract the tens column: 6 – 2 = 4. The results of these

two single-digit subtractions, written side by side, provide the answer:
Subtraction is a bit more complicated if we need to subtract a larger digit from a smaller one.
For example, when subtracting 47 from 92, the units, value (7) of 47 is greater than the units,
value (2) of 92. We can handle this situation using a procedure called borrowing, which is like
carrying in reverse. Ten units can be borrowed from the tens column—that is, from the 9 of
92—leaving 8 in the tens column. Bring the 10 over to the units, column and add it to the 2
already there, giving 12 in that column from which 7 can then be, subtracted:

Complete the subtraction by taking 4 away from 8 in the tens column, which gives 4. The
answer, or difference, is 45.
C. Multiplying positive integers
Multiplication is simply repeated addition and is often, indicated by the times sign (×). The
expression 3 × 4 means that 3 is to be added to itself 4 times or, similarly, that 4 is to be
added to itself 3 times. In either case, the answer is the same: 12. For example, 3 sets of 4
apples together contain a total of 12 apples. When large numbers are involved, however,
such repeated addition is tedious. Multiplication provides a procedure for simplifying
repeated addition. Sometimes a dot or an asterisk is, used instead of a times sign to indicate
the multiplication of two or more numbers, and sometimes parentheses are used. For
example, 3 × 4, 3 · 4, 3 * 4, and (3)(4) all indicate 3 times 4.
The end result of multiplication is called the product. As in addition, it is customary to
memorize the products of the integers between 0 and 9, which are provided in the following
table:
To find the product of any two numbers from 0 to 9, locate one of the numbers in the vertical
column on the left side of the table and the other number in the horizontal row at the top.
The product is the number, in the body of the table that lies at the intersection of the column
and row that have been, selected. For example, 7 × 8 = 56.
In order to multiply numbers with more than one digit, align the units, tens, hundreds, and so
on. The order in which two numbers are, multiplied does not affect their product: 3 × 4 = 12
and 4 × 3 = 12. It is generally easier; however, to multiply two numbers if the number with the
most digits is, placed on the top:

We then multiply each digit of the top number by the bottom number, in this case, 4. Adding
the results of all these multiplications together gives the product 1,544:

Carrying tens and hundreds, as in addition, shortens this operation:

Multiply the 6 by 4, giving 24. Write the 4 in the units place of the product and carry the 2,
which stands for 2 tens, or 20. Multiply the 8 in the tens place by 4, giving 32, then add the
carried 2, giving 34. (We actually multiplied 80 by 4, giving 320, and carried 20, which we
then added to 320.) Write the 4 in the tens place just to the left of the 4 in the units place,
and carry the 3 (which stands for 3 hundreds), placing it over the hundreds column. To
finish, multiply the 3 in the hundreds place by 4, giving 12 (actually 1,200), and add the 3
that we carried, giving 15 (actually 1,500):

We can follow a similar procedure when both numbers to be multiplied have more than

one digit. To multiply 36 by 52, for example, begin by multiplying the top number, 36, by
the unit 2 of the bottom number:
Next multiply the 6 by 5, giving 30, and put the 0 under the number 7 in the tens place of
the partial product. This placement is chosen because the 5 in the bottom number is in the
tens place and actually represents 50. Carry the number 3 as usual. Multiply the 3 in the
tens place of the top number by 5, giving 15, and add the carried 3, giving 18 (really 5 times
30, plus 30, for a total of 180). Now write the 8 in the hundreds place (directly to the left of
the 0 in the tens place), and carry the 1 into the thousands place. We obtain the total
product by adding the two partial products:

For each digit in the bottom number, there will be a line for a partial product obtained as
we multiply out the units, tens, hundreds, and so on. Place each successive partial product
one digit to the left of the partial product above it, corresponding to the position of the
digit being multiply in the bottom number. Thus, in the above example, the partial product
of 5 times 36 yields 180, and the 0 in the digits place must be written in the tens place of
the partial product. This placement results because the 5 in the tens place is really 50 and
therefore the entire partial product is really 1,800:

D. Dividing Positive Integers

The arithmetic operation of division is the opposite, or inverse, of multiplication. Using
the example of 12 divided by 4, we may indicate division by the division sign (12 ÷ 4), a
bar (Ž), a slash (12/4), or the notation p. Division determines how many times one
number is contained in another number. For example, 4 is contained 3 times in 12; thus,
12 apples could be divided into 3 sets of 4 apples, so 12 divided by 4 is 3. The number to
be, divided is, called the dividend, the number the dividend is divided by is called the
divisor, and the end, result of division is called the quotient:
Simple divisions such as 12 ÷ 4 may be, carried out mentally, but more, complicated cases
require a procedure known as long division. Long division involves the repetition of simple
operations. For example, to divide 4,518 by 6, consider the divisor (6) and the first digit (4)
of the dividend to see whether the divisor is contained in that first digit one or more times:

If the first digit is too small (6 is not contained in 4 even once), try to divide the first two
digits of the dividend (45) by the divisor (6) . To determine how many 6s are contained in 45,
make a guess. If we guess 8, we can check our guess by multiplying 6 by 8, which yields 48.
Since 48 is more than 45, the guess was too big. Guessing 6 and multiplying 6 × 6 yields 36—
too small. We know it is too small because when we subtract 36 from 45 (45 – 36) we get 9,
which indicates that 45 contains another 6. Therefore 6 will go into 45 no more than 7 times
(6 × 7 = 42). Write the number 7 in the quotient over the 5 in the dividend, 4,518, and write
the 42 (the product of the divisor, 6, and the first number of the quotient, 7) under the 45
and subtract from it, yielding 3. What we have actually determined so far is that 4,518
contain at least 700 6s, and that 318 is left when these are taken away. In the next step of
the division process, bring the 1 in the dividend down and write it to the right of the 3 to
give 31:

The 6 in the divisor will go into 31 no more than 5 times (6 × 5 = 30). Write the 5 in the

quotient to the right of the 6, above the 1 in the dividend. Place the product of 5 × 6, or 30,
under the 31 and subtract, yielding 1. Bring the 8 from the dividend down and write it to the
right of the 1 to give 18. The 6 in the divisor will go into 18 exactly 3 times, so write the
number 3 in the quotient above the 8 in the dividend:

The answer to how many times 6 will divide 4,518 is therefore 753. We can verify this
solution by multiplying 6 × 753, which yields the dividend 4,518.
 Repetitious operations such as those employed in long division are, called iterative
processes and are often, used in computer science (see Algorithm).

E. Reminders
Dividends are not always evenly divisible by divisors. Had the dividend in the last example
been 4,523 rather than 4,518, the final subtraction would not have yielded 0, but would
have left a remainder of 5—too small for 6 to go into. We can incorporate this remainder
into the answer as a fraction of the divisor, in this case 6, leaving a fractional remainder of
™:

XIV. Multi-digit Division

We can also use long division if a divisor has more than one digit. The method is similar to
the shorter method described above. For example, if we divide 3,626 by 25,

25 will not go into 3 at all, and will go into 36 only once; therefore, we place a 1 over the 6 in
36. Multiply the 25 by 1, write the product (25) under the 36, and subtract from it, yielding
11. Then bring the 2 down to give 112:

Because 25 will go into 112 at most 4 times (4 x 25 = 100), write a 4 next to the 1 in the

quotient, over the 2 in the dividend. Multiply the divisor by 4, write the product (100) under
the 112, and subtract from it, leaving 12. Bring the 6 down for a total of 126. Because 25 will
go into 126 five times, with a remainder of 1, the answer is
XV. Operation with Negative Integers
We can most easily visualize negative numbers by considering the familiar numbers of
arithmetic, the positive integers, arranged in a line and increasing in a positive direction.
Negative numbers are reached by proceeding away from 0 in the opposite direction. The
following number line represents both positive and negative numbers:

To subtract 4 from 2, count backward four units from 2 to reach the difference: -2. (If no sign
appears before a number, we assume the number is positive.)
To handle arithmetic operations involving negative numbers, it is necessary to understand
the concept of absolute value. A number’s absolute value is its magnitude regardless of its
sign. Thus, the absolute value of +5 is 5, and the absolute value of -5 is also 5. In symbolic
notation, the absolute value of any number a is written as |a|. The absolute value of 7, for
example, is written |7| and is equal to 7, while the absolute value of -9 is written |-9| and is
equal to 9.

A.Adding Negative Integers

By referring once again to the number line, which includes both positive and negative
numbers, we can see that the sum of any two negative numbers is negative and equals the
total of the absolute values of the two numbers, but negative:

If, however, the signs of the two numbers are opposite, then we must proceed more
cautiously.

To add -4 and 2, for example, begin at -4 and add 2 by moving 2 units forward, in the positive

direction to the right. The result is -2. The answer is the difference between the absolute
values of the two numbers, with the result taking the sign of the larger number. The answer
to this problem must be negative because in terms of absolute value, |-4| > |2|. (The symbol
> means “greater than.”) We obtain the answer, however, by calculating the difference
between the absolute values 4 and 2. Similarly, to add -3 and 7, begin at -3 and move 7 units
forward, reaching 4. Two rules simplify this process:
Rule 1: To add two numbers with the same sign, add the absolute values of the two
numbers, then, give the answer the sign of the numbers in question:

Rule 2: To add two numbers with different signs, determine their absolute values and

subtract the smaller number from the larger. Give the result the sign of the number with the
larger absolute value:

B.Subtracting Negative Integers

The subtraction of negative numbers is easy to visualize on the number line. Subtracting one
number from another involves moving along the number line in the opposite direction from
addition. To subtract 5 from 7, for example, begin at 7 and move in the negative direction 5
units, reaching 2 as the answer. Similarly, to subtract -2 from 5, start at 5 but reverse
direction and move 2 units in the positive direction to 7. Subtracting -2 basically means
adding 2, or more generally, subtracting a negative number is the same thing as adding the
absolute value of the number. We can therefore change two negative (or minus) signs in a
row into a single positive (or plus) sign. For example:

We can change a negative and a positive sign together into a negative sign:

C.Multiplying Negative Integer

Determining the sign of a product is straightforward. To multiply two numbers with the
same sign, multiply their absolute values and give the resulting product a positive sign:
To multiply two numbers with different signs, multiply their absolute values and give the
resulting product a negative sign:

D. Dividing Negative Integers

Division is the inverse of multiplication; therefore, dividing positive or negative numbers
involves rules similar to those for multiplication. To divide two numbers with the same sign,
divide their absolute values and give the resulting quotient a positive sign:

To divide two numbers with different signs, divide their absolute values and give the resulting
quotient a negative sign:

XVI. Operation with Fraction

Numbers that represent parts of a whole are, called fractions or rational numbers. Simple

fractions are familiar: a dime is s of a dollar; ‚ of a pie plus y of a pie is • of a pie; and so on. In
general, we can express fractions as the quotient of two integers a and b:

The top number in a fraction is, called the numerator and the bottom number is, called the
denominator. Two types of fractions exist: proper and improper. A proper fraction is one in
which the numerator is smaller than the denominator; ’, -, and  are all proper fractions.
An improper fraction is one in which the numerator is larger than the denominator; ”, -ž, and
œ are improper fractions. We can convert improper fractions to mixed fractions or whole
numbers (for example, ” = 1y, -ž = -2, and œ = 2€) by dividing the numerator by the
denominator and expressing any remainder as a fraction of the denominator.
A fraction is said to be reduced to lowest terms if neither the numerator nor the
denominator has a factor in common. A factor is a number by which another number can be,
divided evenly. For example, › is not reduced to lowest terms because both 6 and 8 have 2 as
a factor:

Since a number divided by itself is always equal to 1, = 1. Multiplying any number by 1 does

not change the number, so × • = •. Reduced to lowest terms, then, › is •.

A. Adding and Subtracting Fraction

To add or subtract fractions that have, the same denominators, add or subtract the
numerators according to the rules for integers, and express the result as a fraction of the
denominator. The result is normally, reduced to lowest terms. For example,

Only fractions with equal denominators may be, added or subtracted as they stand. If the

denominators of fractions to be, added are unequal, we must find a common denominator.
In the expression, ’ + •, for example, the denominators 3 and 4 are different. One quick way
to obtain a common denominator for two fractions is to multiply their denominators. In this
case, that multiplication gives us 12. Thus, to add ’ and •, we should change the fractions
into their equivalents with 12 as a common denominator. To do so, multiply both the
numerator and denominator of each fraction by the denominator of the other fraction:
 Multiplying both the numerator and denominator by the same number does not change the
overall value of a fraction. Now we can add

Subtracting fractions requires the same procedure:

B. Multiplying Fraction

Multiplying two fractions, ¡ and ¢, is straightforward. Simply multiply numerators together and

multiply denominators together:

For example, 

The answer has been, reduced to its lowest term, in this case ˜.

The rules for multiplying signed (positive or negative) fractions are the same as those for
multiplying signed integers. The same is true for the rules governing addition, subtraction, and
division of signed fractions:

C. Dividing Fraction
The division of fractions is most easily, understood in terms of reciprocals. Every number
(except 0) has a reciprocal or another number such that the product of the number and its
reciprocal equals one. The reciprocal of 3, for example, is €. Zero (0) has no reciprocal,
because no number can be multiplied by it to equal 1. Any number multiplied by zero equals
zero.
To divide ¡ by ¢, multiply ¡ by the reciprocal of ¢:
 Division is equivalent to multiplying by the reciprocal—that is, y ÷ ‚ is the same as y ×
because both the numerator and the denominator can be multiplied by the same nonzero
number without affecting the overall value of the fraction. Multiply both numerator and
denominator in the graphic above by ¤:

This division is equivalent to multiplying the first number, ¡, by the reciprocal of the second

number—by ¤. The reciprocal of a fraction is simply the fraction flipped upside down. Here is
an example that uses actual numbers:

Dividing a whole number by a fraction works the same way:

XVII. Operation with Decimal

The idea of place values can be, extended to accommodate fractions. Instead of writing 1‘
(one and two-tenths), we can use a decimal point (.) to represent the same fraction as 1.2.
Just as places to the left of the decimal represent units, tens, hundreds, and so on, those to
the right of the decimal represent places for tenths (s), hundredths (t), thousandths ( ),
and so forth. In a decimal number such as 8.632 the numbers to the right of the decimal point
represent

This number is, read “eight and six hundred thirty-two thousandths” or “eight point six three
two.” Zeroes are often, added to the left of the decimal point when a number is less than one.
For example, we can write y as either .5 or 0.5.

A. Adding and Subtracting Decimal

Decimals allow us to add and subtract numbers that include fractions just as we add and
subtract integers. However, we must be careful to always align the decimal point so that tens
are under tens, units under units, tenths under tenths, and so on, ensuring that each value is
being added to or subtracted from a similar value at every step. For example, to add 365.289
and 32.4, align the decimals and then add the numbers beginning at the right and moving to
the left:

Extra zeros to the right of a number do not change the number. Filling in the number of places
with zeroes can help ensure that the same number of places exists to the right of the decimal
point for all numbers being added or subtracted:

The decimal point in the sum falls directly beneath the decimal points in the numbers being
added.
Subtraction with decimals proceeds in much the same way as addition: We can use zeros to
ensure that the numbers’ place values line up, then subtract the numbers as usual.

B. Multiplying Decimal
Multiplying decimals is similar to multiplying integers, except that the position of the decimal
point must be kept in mind. First, multiply decimal numbers as if they were integers, without
considering the decimal points. Then place the decimal point at the appropriate position in the
product so that the number of decimal places is the same as the total number of decimal
places in the numbers being multiplied. For example, in multiplying 0.3 by 0.5

the 15 in the product is the direct value of 3 times 5. We place the decimal point differently
than we do in addition and subtraction. Since the factors, 0.3 and 0.5, each have one decimal
place, the product must have two decimal places. Thus, 0.15 is the product.
Another example will serve to clarify this concept. Multiply 0.2 by 0.3. The product of 2 and 3
is 6, but since 0.2 and 0.3 each have one decimal place, the product must have a total of two
decimal places. We can fulfill this requirement by placing the decimal point two places to the
left of the 6 in the product, then adding a zero to fill the tenths place: 0.06.
More complicated problems are solved similarly:
Because the top number has three decimal places and the bottom number has two, the
product must have a total of five decimal places. Count five places to the left starting with the
digit farthest to the right (the 8 of the final product), adding a zero if necessary. The answer is
0.03108.
One final example will also illustrate the importance of counting the proper number of
decimal places in determining products of decimal numbers. Multiply .001 and .002. The
multiplication of the digits is simple: the answer is 2. However, each number has three
decimal places, giving a total of six places that must be preserved in the product. Insert five
zeroes to fill the places between the 2 and the decimal point. The answer is 0.000002.

C. Dividing Decimal
Like multiplication, the division of decimal numbers follows the same procedures used to
divide integers, except that we must take care to determine the proper placement of decimal
points in quotients. Dividing a decimal number by a whole number is straightforward: Place
the decimal point in the quotient directly above the decimal point in the dividend and ignore
it during the rest of the process of division:

In cases where the divisor is a decimal number, convert the problem to one in which the

divisor is an integer; division may then proceed as in the above example. To divide 14 by 0.7,
for example, convert the divisor to an integer by multiplying it by 10: (0.7)(10) = 7. Then
multiply the dividend by an equal amount. We can understand this procedure more easily by
considering the division rewritten as a fraction. Multiplying both numerator and denominator
by the same amount will not change the value of the fraction:

Similarly, the division of 2.675 by 0.23 can be considered in the form 2.675/0.23. We can

convert this fraction to a division involving an integer divisor, namely 23, if we multiply both
numerator and denominator by 100:
We can convert any division problem involving a decimal divisor into a problem with an
integer divisor simply by moving the decimal point in the divisor as many places to the right as
is necessary to make it an integer. Then move the decimal point in the dividend an equal
number of places to the right, and add zeros if necessary. For example, to divide 21.5 by .002,
move the decimal point in the divisor three places to the right, giving the integer 2. Move the
decimal point in the dividend three places to the right as well:

Carry out the division as usual, placing the decimal point in the quotient directly over the new
decimal point in the dividend. The quotient in this case is 10750:

XVIII. Matrix Theory and Linear Algebra

A. Introduction
Matrix Theory and Linear Algebra, interconnected branches of mathematics that serve as
fundamental tools in pure and applied mathematics and are becoming increasingly
important in the physical, biological, and social sciences.

B.Matrix Theory
A matrix is a rectangular array of numbers or elements of a ring (see Algebra). One of the
principal uses of matrices is in representing systems of equations of the first degree in
several unknowns. Each matrix row represents one equation, and the entries in a row are
the coefficients of the variables in the equations, in some fixed order.

A matrix is usually enclosed in brackets:

In the above matrices, a, b, and c are arbitrary numbers. In place of brackets, parentheses or
double vertical lines may be used to enclose the arrays. The horizontal lines, called rows, are
numbered from the top down; the vertical lines, or columns, are numbered from left to
right; thus, -1 is the element in the second row, third column of M1. A row or column is
called a line.

The size of a matrix is given by the number of rows and columns, so that M1, M2, M3, and M4
are, in that order, of sizes 3 × 3 (3 by 3), 3 × 3, 3 × 2, and 2 × 3. The general matrix of size m ×
n is frequently represented in double-subscript notation, with the first subscript i indicating
the row number, and the second subscript j indicating the column number; a23 is the
element in the second row, third column. This general matrix

may be abbreviated to A = [aij], in which the ranges i = 1, 2, ..., m and j = 1, 2, ..., n should be
explicitly given if they are not implied by the text. If m = n, the matrix is square, and the
number of rows (or columns) is the order of the matrix. Two matrices, A = [aij] and B = [bij],
are equal if and only if they are of the same size and if, for every i and j, aij = bij. The
elements a11, a22, a33, ... constitute the main or principal diagonal of the matrix A = [aij], if it is
square. The transpose AT of a matrix A is the matrix in which the ith row is the ith column of
A and in which the jth column is the jth row of A; thus, from the matrix M3, above,

which is the transpose of M3.

Addition and multiplication of matrices can be defined so that certain sets of matrices form

algebraic systems. Let the elements of the matrices considered be arbitrary real numbers,
although the elements could have been chosen from other fields or rings. A zero matrix is
one in which all the elements are zero; an identity matrix, Im of order m, is a square matrix
of order m in which all the elements are zero except those on the main diagonal, which are
1. The order of an identity matrix may be omitted if implied by the text, and Im is then
shortened to I.

The sum of two matrices is defined only if they are of the same size; if A = [aij] and B = [bij]
are of the same size, then C = A + B is defined as the matrix [cij], in which cij = aij + bij; that is,
two matrices of the same size are added merely by adding corresponding elements. Thus,
in the matrices given above
 The set of all matrices of a fixed size has the property that addition is closed, associative,
and commutative; a unique matrix O exists such that for any matrix A, A + O = O + A = A;
and, corresponding to any matrix A, there exists a unique matrix B such that A + B = B + A =
O.

The product AB of two matrices, A and B, is defined only if the number of columns of the

left factor A is the same as the number of rows of the right factor B; if A = [aij] is of size m ×
n and B = [bjk] is of size n × p, the product AB = C = [cik] is of size m × p, and cik is given by

That is, the element in the ith row and kth column of the product is the sum of the
products of the elements of the ith row of the left factor multiplied by the corresponding
elements of the kth column of the right factor.

C. Linear Algebra
The geometric concept of a vector as a line segment of given length and direction can be
advantageously generalized as follows. An n-vector (n-dimensional vector, vector of order n,
vector of length n) is an ordered set of n elements of a field. As in matrix theory, the
elements are assumed to be real numbers. An n-vector v is represented as: v = [x1, x2, ..., xn]
In particular, the lines of a matrix are vectors; the horizontal lines are row vectors, the
vertical lines are column vectors. The x's are called the components of the vector.

Addition of vectors (of the same length) and scalar multiplication are defined as for
matrices and satisfy the same laws. If w = [y1, y2, ..., yn] and k is a scalar (real number), then
v + w = [x1 + y1, x2 + y2, ..., xn + yn] kv = [kx1, kx2, ..., kxn]

If k1, k2, ..., km are scalars and v1, v2, ..., vm are n-vectors, the n-vector v = k1v1 + k2v2 + ... +
kmvm is called a linear combination of the vectors v1, v2, ..., vm. The m n-vectors are linearly
independent if the only linear combination equal to the zero n-vector, 0 = [0,0, ..., 0], is the
one in which k1 = k2 = ... = km = 0; otherwise, the vectors are linearly dependent. For
example, if v1 = [0, 1, 2, 3], v2 = [1, 2, 3, 4], v3 = [2, 2, 4, 4], v4 = [3, 4, 7, 8], then v1, v2, v3 are
linearly independent, because k1v1 + k2v2 + k3v3 = 0 if and only if k1 = k2 = k3 = 0; v2, v3, and v4
are linearly dependent because v2 + v3 - v4 = 0. If A is a matrix of rank r, then at least one set
of r row, or column, vectors is a linearly independent set, and every set of more than r row,
or column, vectors is a linearly dependent set.

A vector space V is a nonempty set of vectors (see Set Theory), with the properties that (1)
if veV and we V, then v + weV, and (2) if ve V and k is any scalar, then kveV. If S = {vi} is a
set of vectors, all of the same length, all linear combinations of the v's form a vector space
said to be spanned by the v's. If the set B = {w1} spans the same vector space V and is a
linearly independent set, the set B is a basis for V. If a basis for V contains m vectors, every
basis for V will contain exactly m vectors and V is called a vector space of dimension m.
Two- and three-dimensional Euclidean spaces are vector spaces when their points are
regarded as specified by ordered pairs or triples of real numbers. Matrices may be used to
describe linear changes from one vector space into another.

Probability is also theory of probability, branch of mathematics that deals with measuring
or determining quantitatively the likelihood that an event or experiment will have a
particular outcome. Probability is based on the study of permutations and combinations
and is the necessary foundation for statistics.

The foundation of probability is usually ascribed to the 17th-century French mathematicians

Blaise Pascal and Pierre de Fermat, but mathematicians as early as Gerolamo Cardano had
made important contributions to its development. Mathematical probability began in an
attempt to answer certain questions arising in games of chance, such as how many times a
pair of dice must be thrown before the chance that a six will appear is 50-50. Or, in another
example, if two players of equal ability, in a match to be won by the first to win ten games,
are obliged to suspend play when one player has won five games, and the other seven, how
should the stakes be divided?

The probability of an outcome is represented by a number between 0 and 1, inclusive, with

“probability 0” indicating certainty that an event will not occur and “probability 1” indicating
certainty that it will occur. The simplest problems are concerned with the probability of a
specified “favorable” result of an event that has a finite number of equally likely outcomes.
If an event has n equally likely outcomes and f of them are termed favorable, the
probability, p, of a favorable outcome is f/n. For example, a fair die can be cast in six equally
likely ways; therefore, the probability of throwing a 5 or a 6 is 2/6. More involved problems
are concerned with events in which the various possible outcomes are not equally likely. For
example, in finding the probability of throwing a 5 or 6 with a pair of dice, the various
outcomes (2, 3, ... 12) are not all equally likely. Some events may have infinitely many
outcomes, such as the probability that a chord drawn at random in a circle will be longer
than the radius.

Problems involving repeated trials form one of the connections between probability and
statistics. To illustrate, what is the probability that exactly five 3s and at least four 6s will
occur in 50 tosses of a fair die? Or, a person, tossing a fair coin twice, takes a step to the
north, east, south, or west, according to whether the coin falls head, head; head, tail; tail,
head; or tail, tail. What is the probability that at the end of 50 steps the person will be
within 10 steps of the starting point?
In probability problems, two outcomes of an event are mutually exclusive if the probability
of their joint occurrence is zero; two outcomes are independent if the probability of their
joint occurrence is given as the product of the probability of their separate occurrences.
Two outcomes are mutually exclusive if the occurrence of one precludes the occurrence of
the other; two outcomes are independent if the occurrence or nonoccurrence of one does
not alter the probability that the other will or will not occur. Compound probability is the
probability of all outcomes of a certain set occurring jointly; total probability is the
probability that at least one of a certain set of outcomes will occur. Conditional probability is
the probability of an outcome when it is known that some other outcome has occurred or
will occur.

If the probability that an outcome will occur is p, the probability that it will not occur is q = 1
- p. The odds in favor of the occurrence are given by the ratio p:q, and the odds against the
occurrence are given by the ratio q:p. If the probabilities of two mutually exclusive
outcomes X and Y are p and P, respectively, the odds in favor of X and against Y are p to P. If
an event must result in one of the mutually exclusive outcomes O1,O2,..., On, with
probabilities p1,p2,..., pn, respectively, and if v1,v2,...vn are numerical values attached to the
respective outcomes, the expectation of the event is E = p1v1 + p2v2 + ...pnvn. For example, a
person throws a die and wins 40 cents if it falls 1, 2, or 3; 30 cents for 4 or 5; but loses $1.20
if it falls 6. The expectation on a single throw is 3/6 × .40 + 2/6 × .30 - 1/6 × 1.20 = .10.

The most common interpretation of probability is used in statistical analysis. For example,
the probability of throwing a 7 in one throw of two dice is 1/6, and this answer is
interpreted to mean that if two fair dice are randomly thrown a very large number of times,
about one-sixth of the throws will be 7s. This concept is frequently used to statistically
determine the probability of an outcome that cannot readily be tested or is impossible to
obtain. Thus, if long-range statistics show that out of every 100 people between 20 and 30
years of age, 42 will be alive at age 70, the assumption is that a person between those ages
has a 42 percent probability of surviving to the age of 70.

Mathematical probability is widely used in the physical, biological, and social sciences and in
industry and commerce. It is applied in such diverse areas as genetics, quantum mechanics,
and insurance. It also involves deep and important theoretical problems in pure
mathematics and has strong connections with the theory, known as mathematical analysis,
that developed out of calculus

XIX. Set Theory

branch of mathematics, first given formal treatment by the German mathematician Georg
Cantor in the 19th century. The set concept is one of the most basic in mathematics, even
more primitive than the process of counting, and is found, explicitly or implicitly, in every
area of pure and applied mathematics. Explicitly, the principles and terminology of sets are
used to make mathematical statements more clear and precise and to clarify concepts such
as the finite and the infinite.

Set Theory
These diagrams represent ways of grouping objects, or elements, of two sets. R is a subset
of S if every element in R also belongs in S (top left). For example, the odd numbers R =
{1,3,5,7,9} are a subset of the set of whole numbers S = {1,2,3,4,5,6,7,8,9,10}. The union of A
and B consists of all the elements in A and all the elements in B (top right). For example, if A
= {2,3,4,5} and B = {4,5,6,7}, then A È B = {2,3,4,5,6,7}. The intersection of A and B consists of
their common elements (bottom left). For example, A Ç B = {4,5}. The difference between B
and A consists of the elements in B that are not also contained in A (bottom right). For
example, B - A = {6,7}.

A set is an aggregate, class, or collection of objects, which are called the elements of the set.
In symbols, aeS means that the element a belongs to or is contained in the set S, or that the
set S contains the element a. A set S is defined if, given any object a, one and only one of
these statements holds: aeS or aS (that is, a is not contained in S). A set is frequently
designated by the symbol S = { }, with the braces including the elements of S either by writing
all of them in explicitly or by giving a formula, rule, or statement that describes all of them.
Thus, S1 = {2, 4}; S2 = {2, 4, 6, ..., 2n,...} = {all positive even integers}; S3 = {x | x2 - 6x + 11 ≥ 3};
S4 = {all living males named John}. In S3 and S4 it is implied that x is a number; S3 is read as the
set of all xs such that x2 - 6x + 11 ≥ 3.

If every element of a set R also belongs to a set S,R is a subset of S, and S is a superset of R;
in symbols, RÍS, or SÊR. A set is both a subset and a superset of itself. If RÍS, but at least one
element in S is not in R,R is called a proper subset of S, and S is a proper superset of R; in
symbols, RÌS,SÉR. If RÍS and SÍR, that is, if every element of one set is an element of the
other, then R and S are the same, written R = S. Thus, in the examples cited above, S1 is a
proper subset of S2.

If A and B are two subsets of a set S, the elements found in A or in B or in both form a subset
of S called the union of A and B, written AÈB. The elements common to A and B form a
subset of S called the intersection of A and B, written AÇB. If A and B have no elements in
common, the intersection is empty; it is convenient, however, to think of the intersection as
a set, designated by Æ and called the empty, or null, set. Thus, if A = {2, 4, 6}, B ={4, 6, 8, 10},
and C = {10, 14, 16, 26}, then AÈB = {2, 4, 6, 8, 10}, AÈC = {2, 4, 6, 10, 14, 16, 26}, AÇB = {4, 6},
AÇC = Æ. The set of elements that are in A but not in B is called the difference between A
and B, written A - B (sometimes A\B); thus, in the illustration above, A - B ={2}, B - A = {8,
10}. If A is a subset of a set l, the set of elements in l that are not in A, that is, l - A, is called
the complement of A (with respect to l), written l - A = A’ (also written Ā,Ã, ~ A).

The following statements are basic consequences of the above definitions, with A,B,C,...
representing subsets of a set l.

1. AÈB = BÈA.
2. AÇB = BÇA.
3. (AÈB) ÈC = AÈ (BÈC).
4. (AÇB) ÇC = AÇ (BÇC).
5. AÈÆ = A.
6. AÇÆ = Æ.
7. AÈl = l.
8. AÇl = A.
9. AÈ (BÇC) = (AÈB) Ç (AÈC).
10. AÇ (BÈC) = (AÇB) È (AÇC).
11. AÈA’ = l.
12. AÇA’ = Æ.
13. (AÈB)’ = A’ÇB’.
14. (AÇB)’ = A’ÈB’.
15. AÈA = AÇA = A.
16. (A’)’ = A.
17. A - B = AÇB’.
18. (A - B) - C = A - (BÈC).
19. If AÇB = Æ, then (AÈB) - B = A.
20. A - (BÈC) = (A - B) Ç (A - C).

These are laws of the algebra of sets, which is an example of the algebraic system that
mathematicians call Boolean algebra.
If S is a set, the set of all subsets of S is a new set D, sometimes called the derived set of S.
Thus, if S = {a,b,c}; D ={{},{a}, {b},{c}, {a,b}, {a,c}, {b,c}, {a,b,c}. Here,{} is used in place of the
null set Æ, of S; it is an element of D. If S has n elements, the derived set D has 2n elements.
Larger and larger sets are obtained by taking the derived set D2 of D, the derived set D3 of
D2, and so on.

If A and B are two sets, the set of all possible ordered pairs of the form (a,b), with a in A and
b in B, is called the Cartesian product of A and B, frequently written A × B. For example, if A
={1, 2}, B ={x,y,z}, then A × B ={ (1, x), (1, y), (1, z), (2, x), (2, y), (2, z)}. B × A ={ (x, 1), (y, 1), (z,
1), (x, 2), (y, 2), (z, 2)}. Here, A × B≠B × A, because the pair (1, x) must be distinguished from
the pair (x, 1).

The elements of the set A = {1, 2, 3} can be matched or paired with the elements of the set
B = {x,y,z} in several (actually, six) ways such that each element of B is matched with an
element of A, each element of A is matched with an element of B, and different elements of
one set are matched with different elements of the other. For example, the elements may
be matched (1, y), (2, z), (3, x). A matching of this type is called a one-to-one (1-1)
correspondence between the elements of A and B. The elements of the set A = {1, 2, 3}
cannot be put into a 1-1 correspondence with the elements of any one of its proper subsets
and is therefore called a finite set or a set with finite cardinality. The elements of the set B =
{1, 2, 3, ...} can be put into a 1-1 correspondence with the elements of its proper subset C
={3, 4, 5, ...} by matching, for example, n of B with n + 2 of C,n = 1, 2, 3, .... A set with this
property is called an infinite set or a set of infinite cardinality. Two sets having elements that
can be placed in a 1-1 correspondence are said to have the same cardinality.

XX. Number Systems

A. Introduction
Number Systems, in mathematics, various notational systems that have been or are being
used to represent the abstract quantities called numbers. A number system is defined by
the base it uses, the base being the number of different symbols required by the system to
represent any of the infinite series of numbers. Thus, the decimal system in universal use
today (except for computer application) requires ten different symbols, or digits, to
represent numbers and is therefore a base-10 system.
Throughout history, many different number systems have been used; in fact, any whole
number greater than 1 can be used as a base. Some cultures have used systems based on
the numbers 3, 4, or 5. The Babylonians used the sexagesimal system, based on the
number 60, and the Romans used (for some purposes) the duodecimal system, based on
the number 12. The Mayas used the vigesimal system, based on the number 20. The binary
system, based on the number 2, was used by some tribes and, together with the system
based on 8, is used today in computer systems. For historical background, see Numerals.

Number Systems

Binary Hexadecim
Octal Decimal
(Base 2) (Base 8) (Base 10) (Base 16)
0 0 0 0
1 1 1 1
10 2 2 2
11 3 3 3
100 4 4 4
101 5 5 5
110 6 6 6
111 7 7 7
1000 10 8 8
1001 11 9 9
 1010 12 10 A
1011 13 11 B
1100 14 12 C
1101 15 13 D
1110 16 14 E
1111 17 15 F
10000 20 16 10
11111111 377 255 FF
11111010001 3721 2001 7D1

B. Place Value
Except for computer work, the universally adopted system of mathematical notation today
is the decimal system, which, as stated, is a base-10 system. As in other number systems,
the position of a symbol in a base-10 number denotes the value of that symbol in terms of
exponential values of the base. That is, in the decimal system, the quantity represented by
any of the ten symbols used—0, 1, 2, 3, 4, 5, 6, 7, 8, and 9—depends on its position in the
number. Thus, the number 3,098,323 is an abbreviation for (3 × 106) + (0 × 105) + (9 × 104) +
(8 × 103) + (3 × 102) + (2 × 101) + (3 × 100, or 3 × 1). The first “3” (reading from right to left)
represents 3 units; the second “3,” 300 units; and the third “3,” 3 million units. In this
system the zero plays a double role; it represents naught, and it also serves to indicate the
multiples of the base 10: 100, 1000, 10,000, and so on. It is also used to indicate fractions of
integers: 1/10 is written as 0.1, 1/100 as 0.01, 1/1000 as 0.001, and so on.

Two digits—0, 1—suffice to represent a number in the binary system; 6 digits—0, 1, 2, 3, 4,

5—are needed to represent a number in the sexagesimal system; and 12 digits—0, 1, 2, 3, 4,
5, 6, 7, 8, 9, t (ten), e (eleven)—are needed to represent a number in the duodecimal
system. The number 30155 in the sexagesimal system is the number (3 × 6 4) + (0 × 63) + (1 ×
62) + (5 × 61) + (5 × 60) = 3959 in the decimal system; the number 2et in the duodecimal
system is the number (2 × 122) + (11 × 121) + (10 × 120) = 430 in the decimal system

To write a given base-10 number n as a base-b number, divide (in the decimal system) n by

b, divide the quotient by b, the new quotient by b, and so on until the quotient 0 is
obtained. The successive remainders are the digits in the base-b expression for n. For
example, to express 3959 (base 10) in the base 6, one writes
from which, as above, 395910 = 301556. (The base is frequently written in this way as a
subscript of the number.) The larger the base, the more symbols are required, but fewer
digits are needed to express a given number. The number 12 is convenient as a base
because it is exactly divisible by 2, 3, 4, and 6; for this reason, some mathematicians have
advocated adoption of base 12 in place of the base 10.

C. Binary System
The binary system plays an important role in computer technology. The first 20 numbers in
the binary notation are 1, 10, 11, 100, 101, 110, 111, 1000, 1001, 1010, 1011, 1100, 1101,
1110, 1111, 10000, 10001, 10010, 10011, 10100. The zero here also has the role of place
marker, as in the decimal system. Any decimal number can be expressed in the binary
system by the sum of different powers of two. For example, starting from the right,
10101101 represents (1 × 20) + (0 × 21) + (1 × 22) + (1 × 23) + (0 × 24) + (1 × 25) + (0 × 26) + (1 ×
27) = 173. This example can be used for the conversion of binary numbers into decimal
numbers. For the conversion of decimal numbers to binary numbers, the same principle can
be used, but the other way around. Thus, to convert, the highest power of two that does
not exceed the given number is sought first, and a 1 is placed in the corresponding position
in the binary number. For example, the highest power of two in the decimal number 519 is
29 = 512. Thus, a 1 can be inserted as the 10th digit, counted from the right: 1000000000. In
the remainder, 519 - 512 = 7, the highest power of 2 is 2 2 = 4, so the third zero from the
right can be replaced by a 1: 1000000100. The next remainder, 3, consists of the sum of two
powers of 2: 21 + 20, so the first and second zeros from the right are replaced by 1: 519 10 =
10000001112.
Arithmetic operations in the binary system are extremely simple. The basic rules are: 1 + 1 =
10, and 1 × 1 = 1. Zero plays its usual role: 1 × 0 = 0, and 1 + 0 = 1. Addition, subtraction, and
multiplication are done in a fashion similar to that of the decimal system:

Because only two digits (or bits) are involved, the binary system is used in computers, since
any binary number can be represented by, for example, the positions of a series of on-off
switches. The “on” position corresponds to a 1, and the “off” position to a 0. Instead of
switches, magnetized dots on a magnetic tape or disk also can be used to represent binary
numbers: a magnetized dot stands for the digit 1, and the absence of a magnetized dot is the
digit 0. Flip-flops—electronic devices that can only carry two distinct voltages at their outputs
and that can be switched from one state to the other state by an impulse—can also be used
to represent binary numbers; the two voltages correspond to the two digits. Logic circuits in
computers (see Computer; Electronics) carry out the different arithmetic operations of binary
numbers; the conversion of decimal numbers to binary numbers for processing, and of
binary numbers to decimal numbers for the readout, is done electronically.
XXI. Trigonometry
A. Introduction
Trigonometry, branch of mathematics that deals with the relationships between the sides
and angles of triangles and with the properties and applications of the trigonometric
functions of angles. The two branches of trigonometry are plane trigonometry, which deals
with figures lying wholly in a single plane, and spherical trigonometry, which deals with
triangles that are sections of the surface of a sphere.

Using Trigonometry to Find the Height of a Building

To estimate the height, H, of a building, measure the distance, D, from the point of
observation to the base of the building and the angle, θ (theta), shown in the diagram. The
ratio of the height H to the distance D is equal to the trigonometric function tangent θ (H/D =
tan θ). To calculate H, multiply tangent θ by the distance D (H = D tan θ). The angle can be
roughly estimated by pointing one arm at the base of the building and the other arm at the
roof and judging whether the angle formed is close to 15°, 30°, 45°, 60°, or 75°. The angle can
be estimated more accurately with a protractor and a plumb bob made of a pencil hanging
from a string. Hang the plumb bob from the zero point in the middle of the straight edge of
the protractor. Sight along the edge of the protractor at the roof of the building. Measure the
angle formed by the straight edge of the protractor and the plumb bob. Subtract this angle
from 90°.

The earliest applications of trigonometry were in the fields of navigation, surveying, and

astronomy, in which the main problem generally was to determine an inaccessible distance,
such as the distance between the earth and the moon, or of a distance that could not be
measured directly, such as the distance across a large lake. Other applications of
trigonometry are found in physics, chemistry, and almost all branches of engineering,
particularly in the study of periodic phenomena, such as vibration studies of sound, a bridge,
or a building, or the flow of alternating current.
 B. Plane Trigonometry
The concept of the trigonometric angle is basic to the study of trigonometry. A
trigonometric angle is generated by a rotating ray. The rays OA and OB (Fig. 1a, 1b, and 1c)
are considered originally coincident at OA, which is called the initial side. The ray OB then
rotates to a final position called the terminal side. An angle and its measure are considered
positive if they are generated by counterclockwise rotation in the plane, and negative if they
are generated by clockwise rotation. Two trigonometric angles are equal if they are
congruent and if their rotations are in the same direction and of the same magnitude.

An angular unit of measure usually is defined as an angle with a vertex at the center of a

circle and with sides that subtend, or cut off, a certain part of the circumference (Fig. 2).

If the subtended arc s (AB) is equal to one-fourth of the total circumference C, that is, s = ‚C,
so that OA is perpendicular to OB, the angular unit is a right angle. If s = yC, so that the
points A, O, and B are on a straight line, the angular unit is a straight angle. If s = 1/360C, the
angular unit is one degree. If s = C, so that the subtended arc is equal to the radius of the
circle, the angular unit is a radian. By equating the various values of C, it follows that 1
straight angle = 2 right angles = 180 degrees = p radians
 Tangent and Cotangent
Tangent and cotangent are periodic trigonometric functions—that is, their values repeat at
regular intervals called periods. The period of tangent and cotangent is 180°, or p.

Each degree is subdivided into 60 equal parts called minutes, and each minute is subdivided
into 60 equal parts called seconds. For finer measurements, decimal parts of a second may
be used. Radian measurements smaller than a radian are expressed in decimals. The symbol
for degree is °; for minutes, ‘; and for seconds, '. For radian measures either the abbreviation
rad or no symbol at all may be used. Thus

The angular unit radian is understood in the last entry. (The notation 42'.14 may be used
instead of 42.14' to indicate decimal parts of seconds.)

By convention, a trigonometric angle is labeled with the Greek letter theta (θ). If the angle θ
is given in radians, then the formula s = rθ may be used to find the length of the arc s; if θ is
given in degrees, then

C.Trigonometric Function
Trigonometric functions are unitless values that vary with the size of an angle. An angle placed
in a rectangular coordinate plane is said to be in standard position if its vertex coincides with
the origin and its initial side coincides with the positive x-axis.

In Fig. 3, let P, with coordinates x and y, be any point other than the vertex on the terminal
side of the angle θ, and r be the distance between Pand the origin. Each of the coordinates x
and y may be positive or negative, depending on the quadrant in which the point P lies; x may
be zero, if P is on the y- axis, or y may be zero, if P is on the x-axis. The distance r is necessarily
positive and is equal to

in accordance with the Pythagorean theorem (see Geometry).

The six commonly used trigonometric functions are defined as follows:

Since x and y do not change if 2p radians are added to the angle—that is, 360° are added—it
is clear that sin (θ + 2p) = sin θ. Similar statements hold for the five other functions. By
definition, three of these functions are reciprocals of the three others, that is,

If point P, in the definition of the general trigonometric function, is on the y-axis, x is 0;

therefore, because division by zero is inadmissible in mathematics, the tangent and secant of
such angles as 90°, 270°, and -270° do not exist. If P is on the x-axis, y is 0; in this case, the
cotangent and cosecant of such angles as 0°, 180°, and -180° do not exist. All angles have
sines and cosines, because r is never equal to 0.

Cosine and Secant

Cosine and secant are periodic trigonometric functions—that is, their values repeat at regular
intervals called periods. The period of cosine and secant is 360°, or 2p.

Since r is greater than or equal to x or y, the values of sin θ and cos θ range from -1 to +1; tan
θ and cot θ are unlimited, assuming any real value; sec θ and csc θ may be either equal to or
greater than 1, or equal to or less than -1.
 Sine and Cosecant
Sine and cosecant are periodic trigonometric functions—that is, their values repeat at regular
intervals called periods. The period of sine and cosecant is 360°, or 2p.

It is readily shown that the value of a trigonometric function of an angle does not depend on

the particular choice of point P, provided that it is on the terminal side of the angle, because
the ratios depend only on the size of the angle, not on where the point P is located on the
side of the angle.

If θ is one of the acute angles of a right triangle, the definitions of the trigonometric functions

given above can be applied to θ as follows (Fig. 4). Imagine the vertex A is placed at the
intersection of the x-axis and y-axis in Fig. 3, that AC extends along the positive x-axis, and that
B is the point P, so that AB = AP = r. Then sin θ = y/r = a/c, and so on, as follows:

The numerical values of the trigonometric functions of a few angles can be readily obtained;

for example, either acute angle of an isosceles right triangle is 45°, as shown in Fig. 4.
Therefore, it follows that

The numerical values of the trigonometric functions of any angle can be determined

approximately by drawing the angle in standard position with a ruler, compass, and protractor;
by measuring x, y, and r; and then by calculating the appropriate ratios. Actually, it is necessary
to calculate the values of sin θ and cos θ only for a few selected angles, because the values for
other angles and for the other functions may be found by using one or more of the
trigonometric identities that are listed below.

D. Trigonometric identities
The following formulas, called identities, which show the relationships between the
trigonometric functions, hold for all values of the angle θ, or of two angles, θ and φ, for which
the functions involved are:

By repeated use of one or more of the formulas in group V, which are known as reduction

formulas, sin θ and cos θ can be expressed for any value of θ, in terms of the sine and cosine
of angles between 0° and 90°. By use of the formulas in groups I and II, the values of tan θ, cot
θ, sec θ, and csc θ may be found from the values of sin θ and cos θ. It is therefore sufficient to
tabulate the values of sin θ and cos θ for values of θ between 0° and 90°; in practice, to avoid
tedious calculations, the values of the other four functions also have been made available in
tabulations for the same range of θ.

The variation of the values of the trigonometric functions for different angles may be

represented by graphs, as in Fig. 5. It is readily ascertained from these curves that each of
the trigonometric functions is periodic, that is, the value of each is repeated at regular
intervals called periods. The period of all the functions, except the tangent and the
cotangent, is 360°, or 2 p radians. Tangent and cotangent have a period of 180°, or p radians.

Many other trigonometric identities can be derived from the fundamental identities. All are
needed for the applications and further study of trigonometry.
 E. Inverse Function
The statement y is the sine of θ, or y = sin θ is equivalent to the statement θ is an angle, the
sine of which is equal to y, written symbolically as θ = arc sin y = sin-1y. The arc form is
preferred. The inverse functions, arc cos y, arc tan y, arc cot y, arc sec y, arc csc y, are similarly
defined. In the statement y = sin θ, or θ = arc sin y, a given value of y will determine infinitely
many values of θ. Thus, sin 30° = sin 150° = sin (30° + 360°) = sin (150° + 360°). . .= 1/2;
therefore, if θ = arc sin 1/2, then θ = 30° + n360° or θ = 150° + n360°, in which n is any integer,
positive, negative, or zero. The value 30° is designated the basic or principal value of arc sin
1/2. When used in this sense, the term arc generally is written with a capital A. Although
custom is not uniform, the principal value of Arc sin y, Arc cos y, Arc tan y, Arc cot y, Arc sec y,
or Arc csc y commonly is defined to be the angle between 0° and 90° if y is positive; and, if y is
negative, by the inequalities

F. The General Triangle

Practical applications of trigonometry often involve determining distances that cannot be
measured directly. Such a problem may be solved by making the required distance one
side of a triangle, measuring othersides or angles of the triangle, and then applying the
formulas below.

If A, B,  C are the three angles of a triangle, and a, b, c the respective opposite sides, it may be
proved that

The cosine and tangent laws can each be given two other forms by rotating the letters a, b, c
and A, B, C.

These three relationships can be used to solve any triangle, that is, the unknown sides or
angles can be found when one side and two angles, two sides and the included angle, two
sides and an angle opposite one of them (usually there are two triangles in this case), or
when three sides are given.

G. Spherical Trigonometry
Spherical trigonometry, which is used principally in navigation and astronomy, is concerned
with spherical triangles, that is, figures that are arcs of great circles (see Navigation) on the
surface of a sphere. The spherical triangle, like the plane triangle, has six elements, the three
sides a, b, c and the angles A, B, C. But the three sides of the spherical triangle are angular as
well as linear magnitudes, being arcs of great circles on the surface of a sphere and measured
by the angle subtended at the center. The triangle is completely determined when any three
of its six elements are given, since relations exist between the various parts by means of
which unknown elements may be found.

In the right-angled or quadrantal triangle, however, as in the case of the right-angled plane

triangle, only two elements are needed to determine all of the remaining parts. Thus, given c,
A in the right-angled triangle, ABC, with C = 90°, the remaining parts are given by the formula
as sin a = sin c sin A; tan b = tan c cos A; cot B = cos c tan A. When any other two parts are
given the corresponding formulas may be obtained by Napier's rules concerning the relations
of the five circular parts, a, b, complement of c, complement of A, complement of B. With
respect to any particular part, the remaining parts are classified as adjacent and opposite; the
sine of any part is equal to the product of the tangents of the adjacent parts and also to the
product of the cosines of the opposite parts.

In the case of oblique triangles no simple rules have been found, but each case depends on
the appropriate formula. Thus in the oblique triangle ABC, given a, b, and A, the formulas for
the remaining parts are

In spherical trigonometry, as well as in plane, three elements taken at random may not

satisfy the conditions for a triangle, or they may satisfy the conditions for more than one.
The treatment of certain cases in spherical trigonometry is quite formidable, because every
line intersects every other line in two points and multiplies the cases to be considered. The
measurement of spherical polygons may be made to depend upon that of the triangle. If, by
drawing diagonals, one can divide the polygons into triangles, each of which contains three
known or obtainable elements, then all the parts of the polygon can be determined.

Spherical trigonometry is of great importance in the theory of stereographic projection and

in geodesy. It is also the basis of the chief calculations of astronomy; for example, the
solution of the so-called astronomical triangle is involved in finding the latitude and
longitude of a place, the time of day, the position of a star, and various other data.
 H. History
The history of trigonometry goes back to the earliest recorded mathematics in Egypt and
Babylon. The Babylonians established the measurement of angles in degrees, minutes, and
seconds. Not until the time of the Greeks, however, did any considerable amount of
trigonometry exist. In the 2nd century BC the astronomer Hipparchus compiled a
trigonometric table for solving triangles. Starting with 7y° and going up to 180° by steps of
7y°, the table gave for each angle the length of the chord subtending that angle in a circle of
a fixed radius r. Such a table is equivalent to a sine table. The value that Hipparchus used for
r is not certain, but 300 years later the astronomer Ptolemy used r = 60 because the
Hellenistic Greeks had adopted the Babylonian base-60 (sexagesimal) numeration system
(see Mathematics).

In his great astronomical handbook, The Almagest, Ptolemy provided a table of chords for

steps of y°, from 0° to 180°, that is accurate to 1/3600 of a unit. He also explained his method
for constructing his table of chords, and in the course of the book he gave many examples of
how to use the table to find unknown parts of triangles from known parts. Ptolemy provided
what is now known as Menelaus's theorem for solving spherical triangles, as well, and for
several centuries his trigonometry was the primary introduction to the subject for any
astronomer. At perhaps the same time as Ptolemy, however, Indian astronomers had
developed a trigonometric system based on the sine function rather than the chord function
of the Greeks. This sine function, unlike the modern one, was not a ratio but simply the length
of the side opposite the angle in a right triangle of fixed hypotenuse. The Indians used various
values for the hypotenuse.

Late in the 8th century, Muslim astronomers inherited both the Greek and the Indian
traditions, but they seem to have preferred the sine function. By the end of the 10th century
they had completed the sine and the five other functions and had discovered and proved
several basic theorems of trigonometry for both plane and spherical triangles. Several
mathematicians suggested using r = 1 instead of r = 60; this exactly produces the modern
values of the trigonometric functions. The Muslims also introduced the polar triangle for
spherical triangles. All of these discoveries were applied both for astronomical purposes and
as an aid in astronomical time-keeping and in finding the direction of Mecca for the five daily
prayers required by Muslim law. Muslim scientists also produced tables of great precision.
For example, their tables of the sine and tangent, constructed for steps of 1/60 of a degree,
were accurate for better than one part in 700 million. Finally, the great astronomer Nasir al-
Din al-Tusi wrote the Book of the Transversal Figure, which was the first treatment of plane
and spherical trigonometry as independent mathematical sciences.

The Latin West became acquainted with Muslim trigonometry through translations of Arabic

astronomy handbooks, beginning in the 12th century. The first major Western work on the
subject was written by the German astronomer and mathematician Johann Müller, known as
Regiomontanus. In the next century the German astronomer Georges Joachim, known as
Rheticus introduced the modern conception of trigonometric functions as ratios instead of as
the lengths of certain lines. The French mathematician François Viète introduced the polar
triangle into spherical trigonometry, and stated the multiple-angle formulas for sin(nq) and
cos(nq) in terms of the powers of sin(q) and cos(q).

Trigonometric calculations were greatly aided by the Scottish mathematician John Napier,

who invented logarithms early in the 17th century. He also invented some memory aids for
ten laws for solving spherical triangles, and some proportions (called Napier's analogies) for
solving oblique spherical triangles.

Almost exactly one half century after Napier's publication of his logarithms, Isaac Newton

invented the differential and integral calculus. One of the foundations of this work was
Newton's representation of many functions as infinite series in the powers of x (see Sequence
and Series). Thus Newton found the series sin(x) and similar series for cos(x) and tan(x). With
the invention of calculus, the trigonometric functions were taken over into analysis, where
they still play important roles in both pure and applied mathematics.

Finally, in the 18th century the Swiss mathematician Leonhard Euler defined the

trigonometric functions in terms of complex numbers (see Number). This made the whole
subject of trigonometry just one of the many applications of complex numbers, and showed
that the basic laws of trigonometry were simply consequences of the arithmetic of these
numbers.

I. Introduction
Game Theory, mathematical analysis of any situation involving a conflict of interest, with the
intent of indicating the optimal choices that, under given conditions, will lead to a desired
outcome. Although game theory has roots in the study of such well-known amusements as
checkers, tick-tack-toe, and poker—hence the name—it also involves much more serious
conflicts of interest arising in such fields as sociology, economics, and political and military
science.

Aspects of game theory were first explored by the French mathematician Émile Borel, who
wrote several papers on games of chance and theories of play. The acknowledged father of
game theory, however, is the Hungarian-American mathematician John von Neumann, who
in a series of papers in the 1920s and '30s established the mathematical framework for all
subsequent theoretical developments. During World War II military strategists in such areas
as logistics, submarine warfare, and air defense drew on ideas that were directly related to
game theory. Game theory thereafter developed within the context of the social sciences.
Despite such empirically related interests, however, it is essentially a product of
mathematicians.

J. Basic Concepts
In game theory, the term game means a particular sort of conflict in which n of individuals or
groups (known as players) participate. A list of rules stipulates the conditions under which
the game begins, the possible legal “moves” at each stage of play, the total number of moves
constituting the entirety of the game, and the terms of the outcome at the end of play.

a. Move
In game theory, a move is the way in which the game progresses from one stage to another,
beginning with an initial state of the game through the final move. Moves may alternate
between players in a specified fashion or may occur simultaneously. Moves are made either
by personal choice or by chance; in the latter case an object such as a die, instruction card, or
number wheel determines a given move, the probabilities of which are calculable.

b. payoff
Payoff, or outcome, is a game-theory term referring to what happens at the end of a game.
In such games as chess or checkers, payoff may be as simple as declaring a winner or a loser.
In poker or other gambling situations the payoff is usually money; its amount is
predetermined by antes and bets amassed during the course of play, by percentages or by
other fixed amounts calculated on the odds of winning, and so on.

c.Extensive and Normal Form

One of the most important distinctions made in characterizing different forms of games is
that between extensive and normal. A game is said to be in extensive form if it is
characterized by a set of rules that determines the possible moves at each step, indicating
which player is to move, the probabilities at each point if a move is to be made by a chance
determination, and the set of outcomes assigning a particular payoff or result to each
possible conclusion of the game. The assumption is also made that each player has a set of
preferences at each move in anticipation of possible outcomes that will maximize the
player's own payoff or minimize losses. A game in extensive form contains not only a list of
rules governing the activity of each player, but also the preference patterns of each player.
Common parlor games such as checkers and tick-tack-toe and games employing playing
cards such as “go fish” and gin rummy are all examples.
Because of the enormous numbers of strategies involved in even the simplest extensive
games, game theorists have developed so-called normalized forms of games for which
computations can be carried out completely. A game is said to be in normal form if the list of
all expected outcomes or payoffs to each player for every possible combination of strategies
is given for any sequence of choices in the game. This kind of theoretical game could be
played by any neutral observer and does not depend on player choice of strategy.

d. Perfect Information
A game is said to have perfect information if all moves are known to each of the players
involved. Checkers and chess are two examples of games with perfect information; poker
and bridge are games in which players have only partial information at their disposal.

e. Strategy
A strategy is a list of the optimal choices for each player at every stage of a given game. A
strategy, taking into account all possible moves, is a plan that cannot be upset, regardless of
what may occur in the game.

XXII. Kinds of Games

Game theory distinguishes different varieties of games, depending on the number of players
and the circumstances of play in the game itself.

a. One-Person Games
Games such as solitaire are one-person, or singular, games in which no real conflict of interest
exists; the only interest involved is that of the single player. In solitaire only the chance
structure of the shuffled deck and the deal of cards come into play. Single-person games,
although they may be complex and interesting from a probabilistic view, are not rewarding
from a game-theory perspective, for no adversary is making independent strategic choices
with which another must contend.

b. Two-Person Games
Two-person, or dual, games include the largest category of familiar games such as chess,
backgammon, and checkers or two-team games such as bridge. (More complex conflicts—n-
person, or plural, games—include poker, Monopoly, Parcheesi, and any game in which
multiple players or teams are involved.) Two-person games have been extensively, analyzed
by game theorists. A major difficulty that exists, however, in extending the results of two-
person theory to n-person games is predicting the interaction possible among various players.
In most two-party games, the choices and expected payoffs at the end of the game are
generally well-known, but when three or more players are involved, many interesting but
complicating opportunities arise for coalitions, cooperation, and collusion.

c. Zero-Sum Games
A game is, said to be a zero-sum game if the total amount of payoffs at the end of the game is
zero. Thus, in a zero-sum game the total amount won is exactly equal to the amount lost. In
economic contexts, zero-sum games are equivalent to saying that no production or
destruction of goods takes place within the “game economy” in question. Von Neumann and
Oskar Morgenstern showed in 1944 that any n-person non-zero-sum game can be reduced to
an n + 1 zero-sum game, and that such n + 1 person games can be generalized from the
special case of the two-person zero-sum game. Consequently, such games constitute a major
part of mathematical game theory. One of the most important theorems in this field
establishes that the various aspects of maximal-minimal strategy apply to all two-person zero-
sum games. Known as the mini-max theorem, it was first, proven by von Neumann in 1928;
others later succeeded in proving the theorem with a variety of methods in more, general
terms.

XXIII. Applications
Applications of game theory are wide-ranging and account for steadily growing interest in
the subject. Von Neumann and Morgenstern indicated the immediate utility of their work on
mathematical game theory by linking it with economic behavior. Models can be developed,
in fact, for markets of various commodities with differing numbers of buyers and sellers,
fluctuating values of supply and demand, and seasonal and cyclical variations, as well as
significant structural differences in the economies concerned. Here game theory is especially
relevant to the analysis of conflicts of interest in maximizing profits and promoting the
widest distribution of goods and services. Equitable division of property and of inheritance is
another area of legal and economic concern that can be studied with the techniques of game
theory.

In the social sciences, n-person game theory has interesting uses in studying, for example,
the distribution of power in legislative procedures. This problem can be interpreted as a
three-person game at the congressional level involving vetoes of the president and votes of
representatives and senators, analyzed in terms of successful or failed coalitions to pass a
given bill. Problems of majority rule and individual decision making are also amenable to
such study.
Sociologists have developed an entire branch of game theory devoted to the study of issues
involving group decision making. Epidemiologists also make use of game theory, especially
with respect to immunization procedures and methods of testing a vaccine or other
medication. Military strategists turn to game theory to study conflicts of interest resolved
through “battles” where the outcome or payoff of a given war game is either victory or
defeat. Usually, such games are not examples of zero-sum games, for what one player loses
in terms of lives and injuries is not won by the victor. Some uses of game theory in analyses
of political and military events have been criticized as a dehumanizing and potentially
dangerous oversimplification of necessarily complicating factors. Analysis of economic
situations is also usually more complicated than zero-sum games because of the production
of goods and services within the play of a given “game.”

XXIV. Calculus (mathematics)

a. Introduction
Calculus (mathematics), branch of mathematics concerned with the study of such concepts
as the rate of change of one variable quantity with respect to another, the slope of a curve at
a prescribed point, the computation of the maximum and minimum values of functions, and
the calculation of the area bounded by curves. Evolved from algebra, arithmetic, and
geometry, it is the basis of that part of mathematics called analysis.

Calculus (mathematics)

Limits
 This graph, which charts the function f(x)=1x, shows that the value of the function
approaches zero as x becomes larger and larger. Yet even as x approaches infinity, the
value of the function will never quite fall to zero. Zero, therefore, is said to be the limit of
this function.

Calculus (mathematics), branch of mathematics concerned with the study of such concepts as

the rate of change of one variable quantity with respect to another, the slope of a curve at a
prescribed point, the computation of the maximum and minimum values of functions, and the
calculation of the area bounded by curves. Evolved from algebra, arithmetic, and geometry, it
is the basis of that part of mathematics called analysis.

Calculus is widely employed in the physical, biological, and social sciences. It is used, for

example, in the physical sciences to study the speed of a falling body, the rates of change in a
chemical reaction, or the rate of decay of a radioactive material. In the biological sciences a
problem such as the rate of growth of a colony of bacteria as a function of time is easily
solved using calculus. In the social sciences calculus is widely used in the study of statistics
and probability.

Calculus can be applied to many problems involving the notion of extreme amounts, such as

the fastest, the most, the slowest, or the least. These maximum or minimum amounts may
be described as values for which a certain rate of change (increase or decrease) is zero. By
using calculus it is possible to determine how high a projectile will go by finding the point at
which its change of altitude with respect to time, that is, its velocity, is equal to zero. Many
general principles governing the behavior of physical processes are formulated almost
invariably in terms of rates of change. It is also possible, through the insights provided by the
methods of calculus, to resolve such problems in logic as the famous paradoxes posed by the
Greek philosopher Zeno.

The fundamental concept of calculus, which distinguishes it from other branches of

mathematics and is the source from which all its theory and applications are developed, is
the theory of limits of functions of variables (see Function).

Let f be a function of the real variable x, which is denoted f(x), defined on some set of real
numbers surrounding the number x0. It is not required that the function be defined at the
point x0 itself. Let L be a real number. The expression

is read: “The limit of the function f(x), as x approaches x0, is equal to the number L.” The
notation is designed to convey the idea that f(x) can be made as “close” to L as desired
simply by choosing an x sufficiently close to x0. For example, if the function f(x) is defined as
f(x) = x2 + 3x + 2, and if x0 = 3, then from the definition above it is true that
This is because, as x approaches 3 in value, x2 approaches 9, 3x approaches 9, and 2 does
not change, so their sum approaches 9 + 9 + 2, or 20.

Another type of limit important in the study of calculus can be illustrated as follows. Let the
domain of a function f(x) include all of the numbers greater than some fixed number m. L is
said to be the limit of the function f(x) as x becomes positively infinite, if, corresponding to a
given positive number e, no matter how small, there exists a number M such that the
numerical difference between f(x) and L (the absolute value |f(x) - L|) is less than e
whenever x is greater than M. In this case the limit is written as

For example, the function f(x) = 1/x approaches the number 0 as x becomes positively
infinite.

It is important to note that a limit, as just presented, is a two-way, or bilateral, concept:

A dependent variable approaches a limit as an independent variable approaches a
number or becomes infinite. The limit concept can be extended to a variable that is
dependent on several independent variables. The statement “u is an infinitesimal”
meaning “u is a variable approaching 0 as a limit,” found in a few present-day and in
many older texts on calculus, is confusing and should be avoided. Further, it is essential
to distinguish between the limit of f(x) as x approaches x0 and the value of f(x) when x is
x0, that is, the correspondent of x0. For example, if f(x) = sin x/x, then

however, no value of f(x) corresponding to x = 0 exists, because division by 0 is undefined

in mathematics.

The two branches into which elementary calculus is usually divided are differential

calculus, based on the consideration of the limit of a certain ratio, and integral calculus,
based on the consideration of the limit of a certain sum.

b. Differential Calculus
Derivatives
The derivative of a function at a given point is equal to the slope of the line that is
tangent to the function at that given point. In this example, the derivative of f(x) at x0 is
defined as the slope of AB in the limit of h going to zero. As h becomes increasingly
smaller, B moves along the curve towards A, and AB increasingly approximates T, the
tangent to the curve at x0.

Let the dependent variable y be a function of the independent variable x, expressed by y

= f(x). If x0 is a value of x in its domain of definition, then y0 = f(x0) is the corresponding
value of y. Let h and k be real numbers, and let y0 + k = f(x0 + h). (Δx, read “delta x,” is
used quite frequently in place of h.) When Δx is used in place of h,Δy is used in place of k.
Then clearly

and

This ratio is called a difference quotient. Its intuitive meaning can be grasped from the
geometrical interpretation of the graph of y = f(x). Let A and B be the points (x0, y0), (x0 +
h, y0 + k), respectively, as in the Derivatives illustration. Draw the secant AB and the lines
AC and CB, parallel to the x and y axes, respectively, so that h = AC, k = CB. Then the
difference quotient k/h equals the tangent of angle BAC and is therefore, by definition,
the slope of the secant AB. It is evident that if an insect were crawling along the curve
from A to B, the abscissa x would always increase along its path but the ordinate y would
first increase, slow down, then decrease. Thus, y varies with respect to x at different
rates between A and B. If a second insect crawled from A to B along the secant, the
ordinate y would vary at a constant rate, equal to the difference quotient k/h, with
respect to the abscissa x. As the two insects start and end at the same points, the
difference quotient may be regarded as the average rate of change of y = f(x) with
respect to x in the interval AC.

If the limit of the ratio k/h exists as h approaches 0, this limit is called the derivative of y

with respect to x, evaluated at x = x0. For example, let y = x2 and x = 3, so that y = 9. Then
9 + k = (3 + h)2; k = (3 + h)2 - 9 = 6h + h2; k/h = 6 + h; and

Referring back to the Derivatives illustration, the secant AB pivots around A and
approaches a limiting position, the tangent AT, as h approaches 0. The derivative of y
with respect to x, at x = x0, may be interpreted as the slope of the tangent AT, and this
slope is defined as the slope of the curve y = f(x) at x = x0. Further, the derivative of y with
respect to x, at x = x0, may be interpreted as the instantaneous rate of change of y with
respect to x at x0.

If the derivative of y with respect to x is found for all values of x (in its domain) for which
the derivative is defined, a new function is obtained, the derivative of y with respect to x.
If y = f(x), the new function is written as y’ or f’(x), Dxy or Dxf(x), (dy)/(dx) or df(x)/dx.
Thus, if y = x2, y + k = (x + h)2; k = (x + h)2 - x2 = 2xh + h2; k/h = 2x + h, whence

Thus, as before, y’ = f’(x) = 6 at x = 3, or f’(3) = 6; also, f’(2) = 4, f’(0) = 0, and f’(-2) = -4.
Derivatives and Indefinite Integrals of Common Functions

As the derivative f’(x) of a function f(x) of x is itself a function of x its derivative with

respect to x can be found it is called the second (order) derivative of y with respect to x
and is designated by any one of the symbols y” or f”(x), Dx2y or Dx2f(x), (d2y)/(dx2) or
(d2f(x))/(dx2). Third- and higher-order derivatives are similarly designated.

Every application of differential calculus stems directly or indirectly from one or both of

the two interpretations of the derivative as the slope of the tangent to the curve and as
the rate of change of the dependent variable with respect to the independent variable.
In a detailed study of the subject, rules and methods developed by the limit process are
provided for rapid calculation of the derivatives of various functions directly by means of
various known formulas. Differentiation is the name given to the process of finding a
derivative.

Differential calculus provides a method of finding the slope of the tangent to a curve at a

certain point; related rates of change, such as the rate at which the area of a circle
increases (in square feet per minute) in terms of the radius (in feet) and the rate at which
the radius increases (in feet per minute); velocities (rates of change of distance with
 respect to time) and accelerations (rates of change of velocities with respect to time,
therefore represented as second derivatives of distance with respect to time) of points
moving on straight lines or other curves; and absolute and relative maxima and minima.

c.Integral Calculus

Integration
An estimate of the area under a curve can be found by adding the areas of a series of
rectangles whose tops closely match the shape of the curve. As this graph shows, any
rectangle will always have part above the curve and part below the curve. However, if
more rectangles with thinner widths are drawn, there will be less difference between the
parts above and below the curve. To find a precise area from adding the areas of
rectangles, the rectangles would have to be infinitely thin. The calculus technique of
integration provides a way to find the area under a curve.

Let y = f(x) be a function defined for all x’s in the interval [a,b], that is, the set of
x’s from x = a to x = b, including a and b, where a<b (suitable modifications can
be made in the definitions to follow for more restricted ranges or domains). Let
x0, x1, ..., xn be a sequence, ..., hn = xn - xn - 1, in brief, hi = xi - xi - 1, where i = 1, 2, ...,
n. The x’s form a partition of values of x such that a = x0<x1<x2<...<xn - 1 <xn = b,
and let h1 = x1 - x0, h2 = x2 - x1of the interval [a, b]; an h with a value not exceeded
by any other h is called the norm of the partition. Let n values of x, for example,
X1, X2, ..., Xn, be chosen so that xi - 1<Xi<xi, where i = 1, 2, ..., n. The sum of the area
of the rectangles is given by f(X1)h1 + f(X2)h2 + .... + f(Xn)hn usually abbreviated to

(Σ is the Greek capital letter sigma.)

Aside from the given function f(x) and the given a and b, the value of the sum
 clearly depends on n and on the choices of the xi’s and Xi’s. In particular, if,
after the xi’s are chosen, the Xi’s are chosen so that f(Xi), for each i, is a
maximum in the interval [xi - 1, xi] (that is, no ordinate from xi - 1 to xi exceeds the
ordinate at Xi), the sum is called an upper sum; similarly, if, after the xi’s are chosen, the
Xi’s are chosen so that f(Xi), for each i, is a minimum in the interval [xi - 1, xi], the
sum is called a lower sum. It can be proved that the upper and lower sums
will have limits,  and , respectively, as the norm approaches 0. If  and
 are equal and have the common value S, S is called the definite integral of
f(x) from a to b and is written

The symbol ∫ is an elongated S (for sum); the f(x) dx is

suggested by a term f(Xi)hi = f(Xi) Δxi of the sum which is used in defining the
definite integral.

If y = g(x), then by differentiation y’ = g’(x). Let g’(x) = f(x), and C be any

constant. Then f(x) is also the derivative of g(x) + C. The expression g(x) + C is
called the antiderivative of f(x), or the indefinite integral of f(x), and it is
represented by

The dual use of the term integral is justified by one of the fundamental
theorems of calculus, namely, if g(x) is an antiderivative of f(x), then, under
suitable restrictions on f(x) and g(x),

The process of finding either an indefinite or a definite integral of a function f(x) is called
integration; the fundamental theorem relates differentiation and integration.
Integration of a Constant Function
Suppose wire costs 3 cents per centimeter. This constant rate can be expressed as y = f(x)
= 3 for all values of x, where x is the length of wire. The area from 0 to 1 represents the 3
cents that the first centimeter of wire costs. The area from 1 to 2 represents the 3 cents
that the second centimeter of wire costs. All of the area under the “curve” from 0 to 5
represents the total, 3 + 3 + 3 + 3 + 3, that the first five centimeters of wire cost, 15
cents. Adding the area under the curve is represented by the integral symbol ∫. The
integral of the function y = 3 is 3x. To evaluate this integral between the points x=0 and
x=5, subtract the value of 3x at x=0 from the value of 3x at x=5. The total cost is 15 - 0, or
15 cents.

If the antiderivative, g(x), of f(x) is not readily obtainable or is not known, the definite

integral can be approximated by the trapezoidal rule, (b - a) [f(a) + f(b)]/2 or

by the more accurate Simpson’s rule:

If |b - a| is small, Simpson’s rule gives a fairly close result. If |b - a| is large, a good

approximation can be obtained by dividing the interval from a to b into a number of
small intervals and applying Simpson’s rule to the subintervals.

Integral calculus involves the inverse process of finding the derivative of a function, that

is, it is the process of finding the function itself when its derivative is known. For
example, integral calculus makes it possible to find the equation of a curve if the slope of
the tangent is known at an arbitrary point; to find distance in terms of time if the velocity
(or acceleration) is known; and to find the equation of a curve if its curvature is known.
Integral calculus can also be used to find the lengths of curves, the areas of plane and
curved surfaces, volumes of solids of revolution, centroids, moments of inertia, and total
mass and total force.

d. Differential Equation
Calculus leads directly to the branch of mathematics called differential equations, which
is extremely useful in engineering and in the physical sciences. An ordinary differential
equation is an equation involving an independent variable, a dependent variable (one or
both of these two may be missing), and one or more derivatives (at least one derivative
must be present). Many physical laws or statements are initially expressed as differential
equations. For example, the law that the acceleration of gravity is a constant g can be
expressed mathematically by the differential equation d2x/dt2 = g; the principle that the
rate of disintegration of radium is proportional to the amount present is expressed as
dR/dt = -kR. A differential equation is solved if an equivalent equation is found involving
only the independent and dependent variables.

This article has considered functions of a single independent variable only. Partial

derivatives, multiple integrals, and partial differential equations are defined and studied
in investigating functions of two or more independent variables.

e. Development of Calculus
The English and German mathematicians, respectively, Isaac Newton and Gottfried
Wilhelm Leibniz invented calculus in the 17th century, but isolated results about its
fundamental problems had been known for thousands of years. For example, the
Egyptians discovered the rule for the volume of a pyramid as well as an approximation of
the area of a circle. In ancient Greece, Archimedes proved that if c is the circumference
and d the diameter of a circle, then 3‡d<c< 3d. His proof extended the method of
inscribed and circumscribed figures developed by the Greek astronomer and
mathematician Eudoxus. Archimedes used the same technique for his other results on
areas and volumes. Archimedes discovered his results by means of heuristic arguments
involving parallel slices of the figures and the law of the lever. Unfortunately, his treatise
The Method was only rediscovered in the 19th century, so later mathematicians believed
that the Greeks deliberately concealed their secret methods.

During the late middle ages in Europe, mathematicians studied translations of

Archimedes’ treatises from Arabic. At the same time, philosophers were studying
problems of change and the infinite, such as the addition of infinitely many quantities.
Greek thinkers had seen only contradictions there, but medieval thinkers aided
mathematics by making the infinite philosophically respectable.

By the early 17th century, mathematicians had developed methods for finding areas and

volumes of a great variety of figures. In his Geometry by Indivisibles, the Italian
mathematician F. B. Cavalieri, a student of the Italian physicist and astronomer Galileo,
expanded on the work of the German astronomer Johannes Kepler on measuring
volumes. He used what he called “indivisible magnitudes” to investigate areas under the
curves y = xn, n = 1 ...9. Also, his theorem on the volumes of figures contained between
parallel planes (now called Cavalieri’s theorem) was known all over Europe. At about the
same time, the French mathematician René Descartes’La Géométrie appeared. In this
important work, Descartes showed how to use algebra to describe curves and obtain an
algebraic analysis of geometric problems. A codiscoverer of this analytic geometry was
the French mathematician Pierre de Fermat, who also discovered a method of finding
the greatest or least value of some algebraic expressions—a method close to those now
used in differential calculus.

About 20 years later, the English mathematician John Wallis published The Arithmetic of

Infinites, in which he extrapolated from patterns that held for finite processes to get
formulas for infinite processes. His colleague at the University of Cambridge was
Newton’s teacher, the English mathematician Isaac Barrow, who published a book that
stated geometrically the inverse relationship between problems of finding tangents and
areas, a relationship known today as the fundamental theorem of calculus.

Although many other mathematicians of the time came close to discovering calculus, the

real founders were Newton and Leibniz. Newton’s discovery (1665-66) combined infinite
sums (infinite series), the binomial theorem for fractional exponents, and the algebraic
expression of the inverse relation between tangents and areas into methods we know
today as calculus. Newton, however, was reluctant to publish, so Leibniz became
recognized as a codiscoverer because he published his discovery of differential calculus in
1684 and of integral calculus in 1686. It was Leibniz, also, who replaced Newton’s
symbols with those familiar today.

In the following years, one problem that led to new results and concepts was that of
describing mathematically the motion of a vibrating string. Leibniz’s students, the
Bernoulli family of Swiss mathematicians (see Bernoulli, Daniel), used calculus to solve
this and other problems, such as finding the curve of quickest descent connecting two
given points in a vertical plane. In the 18th century, the great Swiss-Russian
mathematician Leonhard Euler, who had studied with Johann Bernoulli, wrote his
Introduction to the Analysis of Infinites, which summarized known results and also
contained much new material, such as a strictly analytic treatment of trigonometric and
exponential functions.
Despite these advances in technique, calculus remained without logical foundations.
Only in 1821 did the French mathematician A. L. Cauchy succeed in giving a secure
foundation to the subject by his theory of limits, a purely arithmetic theory that did not
depend on geometric intuition or infinitesimals. Cauchy then showed how this could be
used to give a logical account of the ideas of continuity, derivatives, integrals, and infinite
series. In the next decade, the Russian mathematician N. I. Lobachevsky and German
mathematician P. G. L. Dirichlet both gave the definition of a function as a
correspondence between two sets of real numbers, and the logical foundations of
calculus were completed by the German mathematician J. W. R. Dedekind in his theory
of real numbers, in 1872.

XXV. Statistics
a. Introduction
Statistics, branch of mathematics that deals with the collection, organization, and
analysis of numerical data and with such problems as experiment design and decision
making.

b. History
The 1990 Census Results By Joseph Gustaitis

It doesn't matter whether it was Horace Greeley who first said, 'Go West, young man, go
West!'—or another newspaperman, John B. L. Soule of the Terre Haute Express.
Americans never needed much urging to head toward the setting sun, and, as the 1990
census clearly demonstrated, they are still on their way west.
By law, the 1990 census's total U.S. population count by state had to be reported by
December 31, 1990, as it was. But that was by no means the end of the process.
Detailed population counts within each state were then steadily reported over the
following weeks to meet an April 1, 1991, deadline, and the debate over the accuracy of
the count went on throughout the year. Nevertheless, the figures showed that between
1980 and 1990, the population of California rose by a robust 25.7 percent, a rate nearly
triple that of the nation as a whole. With a population just under 30 million, California
now has more people than most of the countries in Europe. The fastest-growing state in
the 1980s was neighboring Nevada, where the population spurted a full 50.1 percent, a
rise largely fueled by a boom in the Las Vegas area. The nation's fastest-growing city was
Mesa, Ariz., where the population nearly doubled in just ten years (from 152,404 to
288,091—an 89 percent growth rate). In 1980 the nation's second largest city (New York
being number one) was Chicago; now it is Los Angeles. In all, the 13 states which the
Census Bureau considers the 'West' (including Alaska and Hawaii) had a population
increase in the 1980s of 22.3 percent.
The results of the 1990 census also revealed that the peripatetic American population
has, since the 1950s, amended Soule's advice and has headed south as well as west. The
South Atlantic region, which comprises the coastal states stretching from Delaware to
Florida, witnessed a solid population boost of 17.9 percent, an increase in great
measure driven by mushrooming growth in Florida, where the population went from
under 10 million in 1980 to nearly 13 million in 1990—a rise of 32.7 percent. Florida,
which in 1980 was the seventh largest state, is now fourth. Further to the west, Texas,
which experienced a 19.4 percent growth rate, now contains three of the ten largest
cities in the United States—Houston, Dallas, and San Antonio.
The migration into the Sunbelt and the southwestern deserts makes it feasible to argue
that the individual with the single greatest influence on American population patterns in
the 20th century was Willis Haviland Carrier, the man who invented air conditioning.
Would people have been so eager to flock to these sultry climes without being able to
cool their homes and cars?
The census revealed several other noteworthy things about the U.S. population. It is, for
one, getting older—in 1980, 28 percent of the U.S. population was under 18; by 1990
that figure was down to 26 percent. Meanwhile, the proportion of people 65 and over
rose from 11 to 13 percent. The number of housing units went up 16 percent, while the
population rose only around 10 percent; this is because the average number of persons
living in a U.S. household declined (from 2.75 in 1980 to 2.63 in 1990). In 1980, 60
percent of U.S. households consisted of married couples, with or without children; in
1990 the figure was 55 percent. Women continued to outnumber men; in 1990, 51
percent of the population was female (the same proportion as in 1980). The median
cost of a house went up 5 percent from 1980, to $79,100, and the median rent was
$379, a 19 percent rise from a decade before (there were, of course, wide regional
variations).

The New Congress

The results of the 1990 census will be used for many purposes, from community
planning and business decisions to the distribution of billions of dollars in federal funds.
But the census's primary purpose, under the Constitution, is to provide the data for the
distribution among the states of the 435 seats in the U.S. House of Representatives.
This, in turn, affects the number of votes each state has in the Electoral College (the
number of each state's electors is equal to its number of representatives plus its two
senators). Since the size of the House of Representatives is fixed, states with small
population gains or losses will lose seats to states with greater gains. The first election
held under the reapportionment occurring after the 1990 census will take place in
November 1992, and it is evident that when the 103rd Congress convenes in January
1993, Capitol Hill will host many more suntanned representatives.
As the West and South enjoy population spurts, growth in the Midwest and Northeast
has tapered off. New York, which lost population in the 1970s, experienced a population
growth in the 1980s of 2.5 percent, but this modest hike will not prevent that state from
losing three seats in Congress. This pattern is repeated throughout the Northeast and
Midwest. After the 1980 census these two regions sent 208 representatives to the
House; in 1993 the number will be 193. Florida will pick up four seats in the House;
Texas, three. California will gain no less than seven, for a total of 52, more than any
state has ever had. Moreover, with 54 electoral votes, it alone will have one-fifth of the
total a presidential candidate needs to get elected. As a result, future presidential
hopefuls will be spending a great deal of time and effort in the Golden State.

The Decline of Rural America

Four states lost population during the 1980s—Wyoming, North Dakota, Iowa, and West
Virginia. That all are predominantly rural is notable, because the 1980 census had
detected what was thought at the time to be a significant population upsurge in rural
areas, where the population grew 15 percent in the 1970s, as compared to 11 percent
for the nation. It now appears that the agricultural recession of the 1980s abruptly
turned this much-heralded 'small-town' boom into a bust. Fewer than one out of four
U.S. residents now live in rural areas and small towns, and the shipshape little Main
Streets once idealized by Hollywood increasingly present pictures of boarded-up
storefronts and vacant parking spots. Americans profess nostalgia for small-town life,
but they increasingly end up living elsewhere, specifically in the urban-suburban sprawls
termed by the Census Bureau 'metropolitan areas,' where 90 percent of the nation's
growth in the 1980s took place. The 1990 census figures showed that the United States
has 39 metropolitan areas of at least 1 million people and that just over half of the
national population lives in these 39 areas alone. Together, all of the nation's 284
metropolitan areas contain 77.5 percent of the population.

Problems of Growth
The additional congressional representation that high-growth states will gain gives
them, of course, greater political clout, but an ascending population can be a mixed
blessing. In considering the definition of 'the West,' it is enlightening to heed the remark
of the American historian Bernard De Voto, who said, 'The West begins where the
annual average rainfall drops below 20 inches.' For convenience, he noted, that could be
considered the 100th meridian, a line which runs southward from North Dakota just
east of Bismarck, heads down past Dodge City, Kan., and crosses Texas between
Sweetwater and Abilene. Much of the area west of that line is semi-arid or actual desert.
By early 1991, California was blanching from the effects of a four-year drought and
some citizens, prohibited from watering their lawns, had taken to swabbing the dead
turf with green paint. Though spring rains alleviated the drought, the underlying
problem remained, and heated debate continued over ownership and development of
water rights, along with a movement toward expensive desalination plants. In fast-
growing Las Vegas, which likes to tout itself to tourists as an oasis, demands for water
spurred the Las Vegas Valley Water District to propose drilling 146 wells across Nevada
to tap aquifers beneath the desert—a plan opposed by rural interests and
environmentalists. In the face of such controversies, it remains to be seen whether the
parched American Southwest can support its growth pace with a precarious water
supply.
Florida is not a desert, but there too water is a problem. South Florida's rainfall was 25
to 40 percent below average from 1988 to 1990, and many areas have had to turn to
desalination. Yet even with normal rainfall, the problem remains. The demands of
booming communities for water have placed a heavy strain on Lake Okeechobee and
have siphoned off a good deal of the water supply of the Everglades, which are, to the
dismay of environmentalists, suffering accordingly.
Another concern is Florida's traditional role as a retirement spot. While retirees often
bring considerable money with them, boosting the economy, these same citizens can be
expected to place expensive demands on the state's public health services. Florida
would like to attract more younger families, but to do that it will have to perform an
expensive upgrade of its low-rated educational system. Only time will tell whether
Florida's low tax rates, which lure people just as much as the balmy temperatures, can
remain low under such pressures.

Immigrants and Minorities

On December 26, 1990, the Census Bureau reported that the population of the United
States, as of April 1, 1990, was 249,632,692. This figure included almost a million
overseas military personnel and civilian federal workers who were counted in the 1990
census under a new law; such people were not counted in 1980. The total was 22.2
million more than the 1980 figure, an increase of 10.2 percent. If, as in 1980, overseas
personnel were left out, the increase would be only 9.8 percent. Either way, it was the
smallest percentage increase in any decade in U.S. history, except for the 7.3 percent
recorded for the Depression years of the 1930s.
This slackened growth becomes even more meaningful in light of the fact that between
7 and 9 million immigrants entered the United States in the 1980s. In other words, as
much as 40 percent of the U.S. population growth since the 1980 census was due to an
influx of foreigners—largely from Asia, Latin America, and the Caribbean. The U.S.
fertility rate in the 1980s actually dropped to an all-time historic low (1.86 children per
woman).
Not since the World War I era has the population profile of the United States been so
altered by immigration. While the U.S. population increased by around 10 percent, the
Hispanic population went up over 50 percent, and the Asian population more than
doubled. A quarter of all Americans now describe themselves as nonwhite or of Hispanic
origin, and demographers project that by 2030, over one-third of the U.S. population
will fall into those categories.
The picture of a burgeoning West and South juxtaposed with a slackening Midwest and
Northeast conjures up an image of snowbound families packing up the car and rolling
down the highway toward the sunshine. But to think of U.S. population change in terms
of internal migration would be to miss what may be the most important part of the
picture. A good deal of the growth in many states did not come about as a result of
people relocating from one state to another, but from immigration. Groups traditionally
labeled 'minorities' are becoming majorities in many places. By the year 2000 non-
Hispanic whites could become a minority in California, which now has over one-third of
the Hispanic population of the United States. The Asian population boom in the Golden
State has been equally impressive. In 1980, Asians represented just over 5 percent of
California's population; in 1990 they accounted for nearly 10 percent.
These immigration patterns are repeated throughout the nation. The Asian population
of Texas went up 165.5 percent; the Hispanic population, 45.4 percent. In
Massachusetts both the Asian and Hispanic populations more than doubled. New York,
where the Asian population rose 123.4 percent, now has more people of Asian
background than Hawaii. In short, while most Americans continue to trace their roots to
Europe, their proportion is shrinking dramatically. This change has caused the writer and
demographer Ben Wattenberg to argue that the United States is becoming 'the first
universal nation.'

The Undercount Controversy

The largest minority group in the United States is still the black population, which,
according to the 1990 census, constitutes 12.1 percent of the U.S. total—up from 11.7
percent in 1980. Blacks, however, represent one of the groups most likely to be
undercounted, and many critics found fault with the Census Bureau's tally of the black
population. After the release of the breakdown of the population by race and ethnic
group in March 1991, U.S. Representative Thomas C. Sawyer, the Ohio Democrat who
chairs the House subcommittee that monitors the census, charged that the bureau
missed more than 2 million blacks. Before the count, the Census Bureau estimated the
black population to be 30.6 million. The number subsequently recorded was 29,986,060,
and Sawyer's research led him to propose a figure of 32.4 million. This was of more than
theoretical interest; census figures are used to allocate federal funds, and big-city
mayors complained that the underenumeration seriously harmed their already
precarious financial situation.
Not only blacks are likely to be under-counted; other minority populations, such as
Hispanics and Asians, are too—as well as the homeless. On the night of March 20, 1990,
the Census Bureau sent out workers in an intensive effort to count the homeless, but
the effort was criticized by advocates for the homeless as missing vast numbers. The
Census Bureau found some 230,000 homeless people; estimates of the actual number
range from half a million to as high as 3 million.
The problem of undercounting is illustrated by what happened in Detroit. Preliminary
census figures released in August 1990 estimated that the population of that city had
fallen below 1 million, an important cut-off figure for calculating the city's share of
federal and state dollars. This could have cost Detroit 16 percent of its annual operating
revenue. Volunteers recanvassed the population and turned up over 120,000 people
that the Census Bureau agreed had been left out. As a result, Detroit's population was
placed at 1,027,974, and city officials breathed a sigh of relief.
But if that many people can be overlooked in just one city, how reliable was the 1990
census? In a mobile, diverse population like that of the United States, no one pretends
that the Census Bureau can count everyone. The question is not whether the 1990
census undercounted, but how much. Sawyer contended that the 1990 census reversed
a 50-year trend toward improvement. 'For the first time,' he said, 'instead of getting
better, the undercount is getting worse.'
Some argued that the Census Bureau should adjust the actual census totals upward to
compensate for the estimated undercount, not only in the nation as a whole, but
especially in specific states and localities where the undercount would have the greatest
impact. In 1988 the bureau announced that it would make no statistical adjustments
after the 1990 count, whereupon New York City, New York State, Los Angeles, Chicago,
Houston, and Dade County, Fla. (which includes Miami), filed suit to require such
corrections if they were found necessary. The Commerce Department (the government
agency to which the Census Bureau belongs) then agreed that adjustments might be
made, and the judge left open the possibility that the census results could be contested
if the plaintiffs disagreed with the bureau's final decision on adjustment.
After the 1990 census the Census Bureau conducted a Post Enumeration Survey,
something routinely done after every census. This one covered 165,000 households, and
the results, released in April 1991, indicated that between 4 million and 6 million people
had been missed nationwide. Despite that finding, Secretary of Commerce Robert A.
Mosbacher announced on July 15 that he would not adjust the 1990 census counts. His
argument was that while the actual total count would be rendered more precise by the
adjustment, statistics on how the population is distributed among states, cities, and
neighborhoods might be made less accurate. In short, though it was obvious that the
count was low, it was by no means certain where the missing people were. In making
 this decision, Mosbacher overrode the recommendation of Census Director Barbara
Everitt Bryant that the adjustment be made and raised a furor among politicians whose
constituencies would be most affected, primarily those representing big cities. If the
adjustment had been made, for example, financially hard-pressed New York City would
have stood to gain some 230,000 people—and millions of added dollars in federal aid.
Consequently, the plaintiffs who had filed suit in 1988 did so again, and 12 other
lawsuits were filed to force an adjustment, keep the current count, or at least make the
Commerce Department release the adjusted figures.
The census functions both as a counter of heads and as a spotter of trends. Its
performance of the first task is admittedly flawed, perhaps seriously. But when it comes
to painting a group portrait of the nation, it works in vivid, unmistakable hues. And the
image of America that has emerged is of a restless, dynamically changing country—one
that is becoming more western and southern and one that, perhaps more than ever,
continues to perform its historic role as a haven for immigrants from around the globe.
Source: 1992 Collier’s Year Book.

Domesday Book
Compiled in 1086 under the direction of William the Conquerer, the Domesday Book was
a meticulous survey of feudal estates in England.

Simple forms of statistics have been used since the beginning of civilization, when

pictorial representations or other symbols were used to record numbers of people,
animals, and inanimate objects on skins, slabs, or sticks of wood and the walls of caves.
Before 3000 BC the Babylonians used small clay tablets to record tabulations of
agricultural yields and of commodities bartered or sold. The Egyptians analyzed the
population and material wealth of their country before beginning to build the pyramids
in the 31st century BC. The biblical books of Numbers and 1 Chronicles are primarily
statistical works, the former containing two separate censuses of the Israelites and the
latter describing the material wealth of various Jewish tribes. Similar numerical records
existed in China before 2000 BC. The ancient Greeks held censuses to be used as bases
for taxation as early as 594 BC. See Census.

The Roman Empire was the first government to gather extensive data about the

population, area, and wealth of the territories that it controlled. During the Middle Ages
in Europe few comprehensive censuses were made. The Carolingian kings Pepin the
Short and Charlemagne ordered surveys of ecclesiastical holdings: Pepin in 758 and
Charlemagne in 762. Following the Norman Conquest of England in 1066, William I, king
of England, ordered a census to be taken; the information gathered in this census,
conducted in 1086, was recorded in the Domesday Book. Registration of deaths and
births was begun in England in the early 16th century, and in 1662 the first noteworthy
statistical study of population, Observations on the London Bills of Mortality, was written.
A similar study of mortality made in Breslau, Germany, in 1691 was used by the English
astronomer Edmond Halley as a basis for the earliest mortality table. In the 19th century,
with the application of the scientific method to all phenomena in the natural and social
sciences, investigators recognized the need to reduce information to numerical values to
avoid the ambiguity of verbal description.

At present, statistics is a reliable means of describing accurately the values of economic,

political, social, psychological, biological, and physical data and serves as a tool to
correlate and analyze such data. The work of the statistician is no longer confined to
gathering and tabulating data, but is chiefly a process of interpreting the information.
The development of the theory of probability increased the scope of statistical
applications. Much data can be approximated accurately by certain probability
distributions, and the results of probability distributions can be used in analyzing
statistical data. Probability can be used to test the reliability of statistical inferences and
to indicate the kind and amount of data required for a particular problem.

III STATISTICAL METHODS

How Polls Predict
Professional pollsters typically conduct their surveys among sample populations of 1,000
people. Statistical measurements show that reductions in the margin of error flatten out
considerably after the sample size reaches 1,000.

The raw materials of statistics are sets of numbers obtained from enumerations or

measurements. In collecting statistical data, adequate precautions must be taken to
secure complete and accurate information.

The first problem of the statistician is to determine what and how much data to collect.

Actually, the problem of the census taker in obtaining an accurate and complete count of
the population, like the problem of the physicist who wishes to count the number of
molecule collisions per second in a given volume of gas under given conditions, is to
decide the precise nature of the items to be counted. The statistician faces a complex
problem when, for example, he or she wishes to take a sample poll or straw vote. It is no
simple matter to gauge the size and constitution of the sample that will yield reasonably
accurate predictions concerning the action of the total population.

In protracted studies to establish a physical, biological, or social law, the statistician may

start with one set of data and gradually modify it in light of experience. For example, in
early studies of the growth of populations, future change in size of population was
predicted by calculating the excess of births over deaths in any given period. Population
statisticians soon recognized that rate of increase ultimately depends on the number of
births, regardless of the number of deaths, so they began to calculate future population
growth on the basis of the number of births each year per 1000 population. When
predictions based on this method yielded inaccurate results, statisticians realized that
other limiting factors exist in population growth. Because the number of births possible
depends on the number of women rather than the total population, and because women
bear children during only part of their total lifetime, the basic datum used to calculate
future population size is now the number of live births per 1000 females of childbearing
age. The predictive value of this basic datum can be further refined by combining it with
other data on the percentage of women who remain childless because of choice or
circumstance, sterility, contraception, death before the end of the childbearing period,
and other limiting factors. The excess of births over deaths, therefore, is meaningful only
as an indication of gross population growth over a definite period in the past; the
number of births per 1000 population is meaningful only as an expression of the
proportion of increase during a similar period; and the number of live births per 1000
women of childbearing age is meaningful for predicting future size of populations.

c.Tabulation and Presentation of Data

Frequency-Distribution Table
A frequency-distribution table summarizes data. For example, there were 1200 grades
received on 4 examinations by 10 sections of 30 students each. The first column lists the
ten intervals into which the grades were grouped. The second column lists the midpoints
of these intervals. The third column lists the number of grades in each interval, that is,
their frequency. (There were 20 grades between 0 and 10.) The fourth column lists the
proportion of grades in each interval, that is, their relative frequency. (.017 of the 1200
grades were between 0 and 10.) The fifth column lists the number of grades in an
interval and all intervals below it, that is, their cumulative frequency. (35 grades were in
or below the interval between 10 and 20.) The sixth column lists the proportion of grades
in or below an interval, that is, their relative cumulative frequency. (0.029 of the 1200
grades were in or below the interval 10 to 20.)
The collected data must be arranged, tabulated, and presented to permit ready and
meaningful analysis and interpretation. To study and interpret the examination-grade
distribution in a class of 30 pupils, for instance, the grades are arranged in ascending
order: 30, 35, 43, 52, 61, 65, 65, 65, 68, 70, 72, 72, 73, 75, 75, 76, 77, 78, 78, 80, 83, 85,
88, 88, 90, 91, 96, 97, 100, 100. This progression shows at a glance that the maximum is
100, the minimum 30, and the range, or difference, between the maximum and
minimum is 70.

In a cumulative-frequency graph, such as Fig. 1, the grades are marked on the horizontal
axis and double marked on the vertical axis with the cumulative number of the grades on
the left and the corresponding percentage of the total number on the right. Each dot
represents the accumulated number of students who have attained a particular grade or
less. For example, the dot A corresponds to the second 72; reading on the vertical axis, it
is evident that there are 12, or 40 percent, of the grades equal to or less than 72.

In analyzing the grades received by 10 sections of 30 pupils each on four examinations, a

total of 1200 grades, the amount of data is too large to be exhibited conveniently as in
Fig. 1. The statistician separates the data into suitably chosen groups, or intervals. For
example, ten intervals might be used to tabulate the 1200 grades, as in column (a) of the
accompanying frequency-distribution table; the actual number in an interval, called the
frequency of the interval, is entered in column (c). The numbers that define the interval
range are called the interval boundaries. It is convenient to choose the interval
boundaries so that the interval ranges are equal to each other; the interval midpoints,
half the sum of the interval boundaries, are simple numbers, because they are used in
many calculations. A grade such as 87 will be tallied in the 80-90 interval; a boundary
grade such as 90 may be tallied uniformly throughout the groups in either the lower or
upper intervals. The relative frequency, column (d), is the ratio of the frequency of an
interval to the total count; the relative frequency is multiplied by 100 to obtain the
percent relative frequency. The cumulative frequency, column (e), represents the
number of students receiving grades equal to or less than the range in each succeeding
interval; thus, the number of students with grades of 30 or less is obtained by adding the
frequencies in column (c) for the first three intervals, which total 53. The cumulative
relative frequency, column (f), is the ratio of the cumulative frequency to the total
number of grades.

The data of a frequency-distribution table can be presented graphically in a frequency

histogram, as in Fig. 2, or a cumulative-frequency polygon, as in Fig. 3. The histogram is a
series of rectangles with bases equal to the interval ranges and areas proportional to the
frequencies. The polygon in Fig. 3 is drawn by connecting with straight lines the interval
midpoints of a cumulative frequency histogram.

Newspapers and other printed media frequently present statistical data pictorially by

using different lengths or sizes of various symbols to indicate different values.
 XXVI. Measure of central tendency
After data have been collected and tabulated, analysis begins with the calculation of a
single number, which will summarize or represent all the data. Because data often
exhibit a cluster or central point, this number is called a measure of central tendency.

Let x1, x2, …, xn be the n tabulated (but ungrouped) numbers of some statistic; the most

frequently used measure is the simple arithmetic average, or mean, written , which is
the sum of the numbers divided by n:

If the x's are grouped into k intervals, with midpoints m1, m2, …, mk and frequencies f1, f2,
…, fk, respectively, the simple arithmetic average is given by

with i = 1, 2, …, k.

The median and the mode are two other measures of central tendency. Let the x's be
arranged in numerical order; if n is odd, the median is the middle x; if n is even, the
median is the average of the two middle x's. The mode is the x that occurs most
frequently. If two or more distinct x's occur with equal frequencies, but none with
greater frequency, the set of x's may be said not to have a mode or to be bimodal, with
modes at the two most frequent x's, or trimodal, with modes at the three most frequent
x's.

XXVII. Measure of variability

The investigator frequently is concerned with the variability of the distribution, that is,
whether the measurements are clustered tightly around the mean or spread over the
range. One measure of this variability is the difference between two percentiles, usually the
25th and the 75th percentiles. The pth percentile is a number such that p percent of the
measurements are less than or equal to it; in particular, the 25th and the 75th percentiles
are called the lower and upper quartiles, respectively. The pth percentile is readily found
from the cumulative-frequency graph, (Fig. 1) by running a horizontal line through the p
percent mark on the vertical axis on the graph, then a vertical line from this point on the
graph to the horizontal axis; the abscissa of the intersection is the value of the pth
percentile.
 The standard deviation is a measure of variability that is more convenient than percentile
differences for further investigation and analysis of statistical data. The standard
deviation of a set of measurements x1, x2, …, xn, with the mean  is defined as the
square root of the mean of the squares of the deviations; it is usually designated by the
Greek letter sigma (σ). In symbols

The square, σ2, of the standard deviation is called the variance. If the standard deviation
is small, the measurements are tightly clustered around the mean; if it is large, they are
widely scattered.

XXVIII. Correlation
When two social, physical, or biological phenomena increase or decrease proportionately
and simultaneously because of identical external factors, the phenomena are correlated
positively; under the same conditions, if one increases in the same proportion that the
other decreases, the two phenomena are negatively correlated. Investigators calculate
the degree of correlation by applying a coefficient of correlation to data concerning the
two phenomena. The most common correlation coefficient is expressed as

in which x is the deviation of one variable from its mean, y is the deviation of the other
variable from its mean, and N is the total number of cases in the series. A perfect positive
correlation between the two variables results in a coefficient of +1, a perfect negative
correlation in a coefficient of -1, and a total absence of correlation in a coefficient of 0.
Intermediate values between +1 and 0 or -1 are interpreted by degree of correlation.
Thus, .89 indicates high positive correlation, -.76 high negative correlation, and .13 low
positive correlation.

XXIX. Mathematical model

Distribution of IQ Scores
The distribution of scores (commonly called IQ scores) on the Wechsler Adult Intelligence
Scale follows an approximately normal curve, an average distribution of values. The test
is regularly adjusted so that the median score is 100—that is, so that half of the scores
fall above 100, and half fall below.

A mathematical model is a mathematical idealization in the form of a system,

proposition, formula, or equation of a physical, biological, or social phenomenon. Thus, a
theoretical, perfectly balanced die that can be tossed in a purely random fashion is a
mathematical model for an actual physical die. The probability that in n throws of a
mathematical die a throw of 6 will occur k times is

in which (À) is the symbol for the binomial coefficient

The statistician confronted with a real physical die will devise an experiment, such as
tossing the die n times repeatedly, for a total of Nn tosses, and then determine from the
observed throws the likelihood that the die is balanced and that it was thrown in a
random way.

In a related but more involved example of a mathematical model, many sets of

measurements have been found to have the same type of frequency distribution. For
example, let x1, x2, …, xN be the number of 6's cast in the N respective runs of n tosses of
 a die and assume N to be moderately large. Let y1, y2, …, yN be the weights, correct to the
nearest 1/100 g, of N lima beans chosen haphazardly from a 100-kg bag of lima beans.
Let z1, z2, …, zN be the barometric pressures recorded to the nearest 1/1000 cm by N
students in succession, reading the same barometer. It will be observed that the x's, y's,
and z's have amazingly similar frequency patterns. The statistician adopts a model that is
a mathematical prototype or idealization of all these patterns or distributions. One form
of the mathematical model is an equation for the frequency distribution, in which N is
assumed to be infinite:

in which e (approximately 2.7) is the base for natural logarithms (see Logarithm). The
graph of this equation (Fig. 4) is the bell-shaped curve called the normal, or Gaussian,
probability curve. If a variate x is normally distributed, the probability that its value lies
between a and b is given by

The mean of the x's is 0, and the standard deviation is 1. In practice, if N is large, the
error is exceedingly small.

XXX. Tests of reliability

 The statistician is often called upon to decide whether an assumed hypothesis for some
phenomenon is valid or not. The assumed hypothesis leads to a mathematical model; the
model, in turn, yields certain predicted or expected values, for example, 10, 15, 25. The
corresponding actually observed values are 12, 16, 21. To determine whether the
hypothesis is to be kept or rejected, these deviations must be judged as normal
fluctuations caused by sampling techniques or as significant discrepancies. Statisticians
have devised several tests for the significance or reliability of data. One is the chi-square
(c2) test. The deviations (observed values minus expected values) are squared, divided by
the expected values, and summed:

The value of c2 is then compared with values in a statistical table to determine the
significance of the deviations.

XXXI. higher statistics

The statistical methods described above are the simpler, more commonly used methods
in the physical, biological, and social sciences. More advanced methods, often involving
advanced mathematics, are used in further statistical studies, such as sampling theory,
inference and estimation theory, and design of experiments.