AMM - Vol.118 NR 08 PDF

THE AMERICAN MATHEMATICAL
MONTHLY
VOLUME 118, NO. 9 NOVEMBER 2011
The Chief Accountant and Mathematical Friend of 767

Ramanujan—S. Narayana Aiyar
Bruce C. Berndt
On Forming Committees 777

P. Haxell
Complete Dissections: Converting Regions and Their 789

Boundaries
Tom M. Apostol and Mamikon A. Mnatsakanian
A “Bouquet” of Discontinuous Functions for Beginners in 799

Mathematical Analysis
Giacomo Drago, Pier Domenico Lamberti, and Paolo Toni
An Introduction to Gauss Factorials 812

John B. Cosgrave and Karl Dilcher
NOTES
Enumerating the Rationals from Left to Right 830

S. P. Glasby
On the Banach-Steinhaus Theorem for Topological Groups 836

Semyon N. Litvinov
Ramanujan’s “Most Beautiful Identity” 839

Michael D. Hirschhorn
PROBLEMS AND SOLUTIONS 846
REVIEWS
Elementary Probability for Applications. By Rick Durrett 854

Matthew Richey
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THE AMERICAN MATHEMATICAL
MONTHLY
VOLUME 118, NO. 9 NOVEMBER 2011
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The Chief Accountant and Mathematical
Friend of Ramanujan—S. Narayana Aiyar
Bruce C. Berndt
Abstract. S. Narayana Aiyar was chief accountant at the Madras Port Trust office, where Ra-
manujan worked as a clerk in 1912–1914. In this article, a short biography of S. Narayana
Aiyar is given, his mathematical contributions are discussed, and his personal and mathemat-
ical relationships with Ramanujan are examined and emphasized.
1. INTRODUCING S. NARYANA AIYAR. In June 2009, the author had the priv-
ilege of participating in a conference commemorating K. Venkatachaliengar on the
centenary of his birth. For many years, more than any other Indian mathematician,
Venkatachaliengar had kept the spirit of Ramanujan’s mathematics alive in his home-
land, and so it was fitting that a conference in his memory be held on the beautiful
campus of Infosys in Bangalore. At this five-day meeting, a film entitled “God, Zero
and Infinity” about the life of Ramanujan was screened for the delegates. At one partic-
ular juncture in the film, I was saddened to hear one of the narrators condescendingly
proclaim that S. Narayana Aiyar, chief accountant at the Madras Port Trust office, was
an “amateur mathematician,” who was unable to appreciate the work of his employee,
S. Ramanujan. Since S. Narayana Aiyar was not able to defend himself from this un-
charitable remark, this listener decided to construct a defense for him.
Readers familiar with the life of Ramanujan might recall that after receiving a
monthly stipend from R. Ramachandra Rao for more than one year, in February, 1912,
Ramanujan wrote the Madras Port Trust office seeking employment as a clerk. Having
obtained a position at the Port Trust office, Ramanujan was supervised by S. Narayana
Aiyar. More than any other person in India during the years 1910–1914 when Ramanu-
jan lived in Madras prior to his departure for England in March, 1914, S. Narayana
Aiyar was able to appreciate Ramanujan’s work. It has been related that during the
evenings and into the wee hours of the mornings, S. Narayana Aiyar and Ramanujan
worked together on mathematics. Of course, we are not claiming that S. Narayana Ai-
yar was anywhere near the equal of Ramanujan in mathematical ability, but clearly
Ramanujan had respect for him, both as chief accountant at the Madras Port Trust
office and as a mathematician. Since many readers will have some knowledge and
interest in Ramanujan and his work, but will have little familiarity with S. Narayana
Aiyar and his influence on Ramanujan, in this brief article we provide a few details
about the life of S. Narayana Aiyar. Some of our narration can be found in the author’s
article in [9, pp. 97–98]. However, in contrast to our account in [9], we also discuss his
mathematical contributions.
Indian names when translated into English often bear different spellings. Another
rendition of Aiyar is Iyer, which is the preferred family spelling. Since Aiyar was used
in Narayana Aiyar’s mathematical publications, we adhere to this spelling.
2. CULTURAL BACKGROUND. We first discuss the religious and cultural back-

ground in which Ramanujan and Narayana Aiyar lived to help us understand the re-
lationship between them. Ramanujan and Narayana Aiyar were Tamils, and in Tamil-
http://dx.doi.org/10.4169/amer.math.monthly.118.09.767
November 2011] S. NARAYANA AIYAR—FRIEND OF RAMANUJAN 767

speaking areas of South India, in particular, in the southern Indian state of Tamil Nadu
(where Madras, now Chennai, is located), the two primary Hindu Brahmin groups are
the Siva-worshiping Iyers (also called Aiyars) and the Vishnu-worshiping Iyengars.
Thus, although both were Brahmins, Ramanujan was an Iyengar, and Narayana Aiyar
was an Iyer. Orthodox Iyengars are generally more conservative than orthodox Iyers,
but perhaps more so at that time. Narayana is another name for the god Vishnu, but
the Iyers also use that name! For generations, this community were priests, who were
involved mainly in the study of Hindu Sanskrit scriptures and the Vedas, while observ-
ing a strict personal and community code. The Vedas (meaning knowledge or wisdom)
are the oldest Hindu scriptures. Nonmaterialistic, living simply, and fully backed by
the early Hindu rulers, these two Brahmin conclaves were in the top hierarchy of the
Hindu caste system, which was a gradation based on professions. The British found
it convenient to tap this readily available resource for the administration of the vast
Indian continent. Consequently, the Brahmin community gradually began to deviate
from their original religious and social pursuits, and so from the mid-1800s to the
mid-1900s, there was an explosion from the community of legal luminaries, adminis-
trators, and mathematicians that included the prestigious Indian Civil Service, social
reformers, and, unfortunately for the British, the Indian freedom fighters. This is a
brief history of the background and the social climate in which Narayana Aiyar and
Ramanujan were born and lived.
3. LIFE STORY OF S. NARAYANA AIYAR. S. Narayana Aiyar was born on De-

cember 15, 1874 in Cumbum, near Madurai in Tamil Nadu. (The initial S. stands for
his father’s name Subbanarayanan.) His parents, Subbanarayanan Iyer and Lakshmi,
were poor. Both his father and paternal grandfather, Ananthanarayana Iyer, were Hindu
priests, well versed in Sanskrit and the Vedas, and made their livings conducting Hindu
rituals in private homes. Narayana Aiyar’s family was so poor that once during his
early school days, at the annual school inspection, he had to share his only shirt with
his elder brother. Having struggled in his early life, he held a lifelong affectionate con-
cern for the welfare of the poor and everyone with whom he worked. He had one elder
brother and three sisters.
Narayana Aiyar earned an M.A. in mathematics from St. Joseph’s College in
Trichinopoly, one of the premier colleges in southern India at that time. In those
days, the M.A. degree was the highest degree that was offered in any subject in the
absence of the Ph.D., which was introduced only many years later. While teaching at
St. Joseph’s College, Narayana Aiyar became known to Francis Spring, who at that
time was serving with the railways in Trichinopoly. In 1900, Narayana Aiyar was
asked to become office manager at the Madras Port Trust office, at the invitation of
Francis Spring, who by that time had become chairman of the Madras Port Trust, who
was to be knighted in 1911, and who was to become one of Ramanujan’s most ardent
supporters. In due course, Narayana Aiyar was promoted to chief accountant, his po-
sition at the time of Ramanujan’s application for employment as a clerk in February,
1912, and thus became the highest-ranking Indian at the Port Trust, where he served
until his retirement in 1934. In recognition of his service, the government awarded
him the title of “Rao Bahadur,” a title of honor issued by the British to individuals who
had performed great service to the nation.
Narayana Aiyar and his wife Avibakthanayaki had two daughters, Meenakshi and
Kamakshi, and one son, Subbanarayanan. In 1925, Narayana Aiyar’s son-in-law (hus-
band of Kamakshi), M. S. Venkataraman (1902–1967), was appointed to fill a vacancy
in the Madras Port Trust office. At the time of Indian independence in 1947, he was
768 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

S. Narayana Aiyar (1874–1937)
traffic manager and then became the first Indian to officiate as chairman of the Madras
Port Trust. Upon assuming the post, he gathered together all of the files connected with
Ramanujan and transferred them to the National Archives in New Delhi. Copies of
these papers were secured for the present author by Nobel laureate S. Chandrasekhar,
and many are reproduced in the author’s book Ramanujan: Letters and Commen-
tary, coauthored with R. A. Rankin [7]. Narayana Aiyar’s son, Subbanarayanan, also
joined the Madras Port Trust, serving as assistant section master and retiring as office
manager.
One incident in Narayana Aiyar’s life is often related with amusement by his de-
scendants. Narayana Aiyar’s wife had arranged a prospective alliance for her younger
sister and requested that Narayana Aiyar visit the young man and give his valued opin-
ion. Accordingly, Narayana Aiyar called on the prospective bridegroom and asked him
but one question. Satisfied with the reply, he gave his approval for the marriage. The
question he asked was, “What is the value of π ?”
Narayana Aiyar was known for his simplicity bordering on austerity. Although he
could have afforded a more comfortable mode of transportation, he traveled to and
from work by tramcar, the ordinary transport for common people. He did not even
own a radio, although in the evenings he often visited his daughter Kamakshi to listen
to the English news on her radio. He insisted on absolute silence while taking food, as
prescribed in the scriptures. He was religious and scrupulously honest, had an indepen-
dent mind, and was a (silent) champion of Indian independence. Although he suffered
from a mild case of diabetes, he took good care of his health and lived a normal life
until shortly before he died on January 17, 1937.
As indicated above, after working at the Madras Port Trust office during the day,
Narayana Aiyar and Ramanujan often worked on mathematics together at night. It is
interesting that Narayana Aiyar communicated to the Journal of the Indian Mathe-
matical Society two brief papers summarizing some of Ramanujan’s discoveries in the
theory of prime numbers, and some applications of his “master theorem” for evaluat-
ing integrals [3, 4]. By the time these two papers were communicated, Ramanujan had
already published three papers and several problems in the Journal of the Indian Math-
ematical Society, and so one might naturally ask why Ramanujan did not communicate
these results himself to the Journal of the Indian Mathematical Society. The answer

to this question can be implicitly discerned from the initial correspondence between
Ramanujan and Hardy. Hardy had eagerly requested some of Ramanujan’s proofs, es-
pecially in the theory of primes. However, Ramanujan was very reluctant to provide
such proofs. We think that the reason for this was that Ramanujan was well aware
that his proofs in prime number theory were not rigorous. In fact, later Ramanujan did
send to Hardy some of his arguments in the theory of prime numbers, and indeed in his
reply, Hardy carefully pointed out to Ramanujan that certain processes were not jus-
tified [7, pp. 87–89]. Similarly, Ramanujan’s powerful procedure for evaluating large
classes of integrals and summing many infinite series was lacking in rigor. Because
Ramanujan was likely afraid to submit papers for publication wherein his results were
not rigorously proven, requesting Narayana Aiyar to submit research announcements
was a means by which he could claim some credit without subjecting his methodology
to the demands of rigor.
Ramanujan’s first two letters to Hardy were composed with the help of Narayana
Aiyar, who strongly advised Ramanujan to accept Hardy’s invitation to come to Cam-
bridge. Narayana Aiyar and Ramanujan journeyed to the temple at Namakkal, where
Goddess Namagiri gave Ramanujan permission to accept Hardy’s bidding. Narayana
Aiyar then had to arrange with Sir Francis Spring concessions in order that Ramanujan
might take leave of his position with the Port Trust. Although Sir Francis Spring un-
failingly supported Ramanujan, he adhered to strict rules in the office. Narayana Aiyar
often jokingly referred to him at home as “oosi milagai,” a Tamil name for a variety of
hot green chilies.
4. MATHEMATICAL CONTRIBUTIONS. In 1907, Narayana Aiyar became a

founding member of the Indian Mathematical Club, which soon afterward became
known as the Indian Mathematical Society. He served as the first assistant secre-
tary from 1907 to 1910, and later served as treasurer from 1914 to 1928. Clearly,
administrative duties in the positions he held after he left St. Joseph’s College fore-
stalled mathematical activity, but in 1907 he evidently gained new life, for he began
to contribute problems and solutions to the Educational Times. Two years later, he
commenced to submit and solve problems in the newly founded Journal of the Indian
Mathematical Club, which in 1912 became the Journal of the Indian Mathematical
Society. In addition to the two communications of Ramanujan’s theorems to the Indian
Mathematical Society, Narayana Aiyar himself published two further papers [5, 6] of
his own in the Journal of the Indian Mathematical Society and an earlier one in the
Journal of the Indian Mathematical Club [2].
Early in the twentieth century, mathematics taught in Indian schools and colleges
reflected what was taught in England. The emphasis was on solving problems, in-
stead of developing theories. The content of articles published in the Journal of the
Indian Mathematical Society was not dissimilar in spirit to the questions and solutions
published in the Journal. Moreover, if one examines Ramanujan’s early papers and
problems published in the Journal of the Indian Mathematical Society, several clearly
reflect the nature of mathematics published in the same issues of the Journal and taught
in India at that time. And Narayana Aiyar’s problems and solutions also mirror the kind
of mathematics popular in then-contemporary India. Narayana Aiyar’s problems and
solutions in the Educational Times will first be described. Then we briefly discuss his
questions and solutions in the Journal of the Indian Mathematical Club and the Jour-
nal of the Indian Mathematical Society. The average difficulty of problems submitted
to the Educational Times is perhaps higher than for problems submitted to the Journal
of the Indian Mathematical Club and the Journal of the Indian Mathematical Society,
and so we concentrate on the former set of problems.
770 c THE MATHEMATICAL ASSOCIATION

OF AMERICA [Monthly 118
During the years 1907–1911, S. Narayana Aiyar submitted a total of 26 problems to
the Educational Times. During the same period, 14 of his solutions to problems posed
by other contributors were also published in the Educational Times. In fact, except
for one contributed solution, his proposed questions and solutions were all published
during the years 1907–1911. Of the 26 posed problems, 12 are in analysis, mostly on
integrals; 13 are in geometry, with most appearing in 1909 and 1910; and the remain-
ing question was on determinants. Of the problems solved by Narayana Aiyar, only
three were in analysis; four were in algebra, including determinants; seven were in ge-
ometry, toward the end of his period of solutions; one was in combinatorics; and one
was in probability. Many of Narayana Aiyar’s problems, particularly in analysis and
geometry, were extraordinarily difficult. In fact, for 13 of the 26 problems that he sub-
mitted, no solutions were published. Most of Narayana Aiyar’s problems in geometry
are on elementary Euclidean geometry (mainly on conics), and were solved syntheti-
cally in the Apollonius style. The term “synthetic geometry” refers to Euclidian-type
axiomatic geometry, in contrast to organizations of geometries based on coordinates
or other analytic means. In personal communications, Illinois colleague John Wetzel
opined, “I think it takes considerable ingenuity to use synthetic methods to solve such
questions,” while Greek geometer A. Hatzipolakis remarked “they are interesting, and
in my opinion deserve the attention of geometers.”
Those submitting and solving problems in the Educational Times, in addition to
S. Narayana Aiyar, included several Indian mathematicians who also prominently con-
tributed articles and problems to the Journal of the Indian Mathematical Society. These
included M. T. Naraniengar, S. Narayanan, K. J. Sanjana, and V. Ramaswami Aiyar,
the founder of the Indian Mathematical Society and the one whom Ramanujan ap-
proached in 1910 seeking a position in his office. Narayana Aiyar solved three prob-
lems submitted by Sanjana, who solved one of Narayana Aiyar’s problems. Recall that
the second published paper by Ramanujan is entitled “On question 330 of Professor
Sanjana” [12], [13, pp. 15–17]. “Question 330” refers to a problem posed by Sanjana
in the problems section of the Journal of the Indian Mathematical Society. In [12],
Ramanujan cleverly summed a class of infinite series, generalizing that in Sanjana’s
problem. A problem posed by V. Ramaswami Aiyar was the one in elementary proba-
bility solved by S. Narayana Aiyar.
Although known everywhere by its abbreviated title, The Educational Times, the
official name of this monthly periodical devoted to educational themes was the Educa-
tional Times and Journal of the College of Preceptors. The Educational Times was a
well-respected periodical that was published in London and known throughout Europe
and, of course, in India. It was first published in 1847 under the title The Educational
Times, and it has continued up to the present day under several names with the current
one being The Educational Times: A Digest of Current Educational Literature.
As indicated above, each issue contained a list of newly published problems as well
as solutions to those previously published. At the end of the year, a separate volume
entitled Mathematical Questions and Solutions from the “Educational Times” was
published; herein we find the complete solutions that were published during the pre-
vious year. These volumes also contained occasional articles, and S. Narayana Aiyar
published one such article in 1907 [1].
G. H. Hardy published 28 original problems and 16 solutions in The Educational
Times. (Ironically, these numbers are almost identical with those of Narayana Aiyar.)
The first publication of Hardy of any kind was Problem 13848, which appeared
on May 1, 1898, and which asked for the general equation, in symmetric form, of
a circle through two fixed points. His solution appeared in the issue of August 1,
1898. (Note that the lag time between the publication of a problem and its solution

was considerably shorter in 1898 than it is with problem sections in contemporary
journals.)
The first problem (#11042) that Narayana Aiyar solved was one that had been sub-
mitted in June, 1891, but evidently solutions were not received until late 1907 or early
1908 when three solutions, including that of Narayana Aiyar, were published in the
February 1, 1908 issue. The problem is on a recurrence relation for Legendre polyno-
mials. It is interesting to note that the second published solution was by G. N. Watson,
who wrote over 30 papers on Ramanujan’s work in the late 1920s and 1930s, and who,
up until recent times, was more acquainted with Ramanujan’s work than any other
mathematician.
S. Narayana Aiyar’s first question (#16219) was published in the issue of June 1,
1907 and is of historical interest. Narayana Aiyar had discovered that the solutions of
certain differentials given by A. Cayley in his treatise An Elementary Treatise on Ellip-
tic Functions were incorrect [10, Art. 414]. Narayana Aiyar’s problem was to find the
correct values. (Before we proceed further, it might be remarked that few “amateur”
mathematicians (if any) would have been capable in those (or any) years of reading
Cayley’s Elliptic Functions.) Over the years, there has been extensive debate on the
sources (if any) from which Ramanujan learned about elliptic functions before depart-
ing for England in 1914. And, if he had seen texts on the subject, was he aware of
Cayley’s treatise? There is considerable, although slightly conflicting, evidence that
Ramanujan had been familiar with Cayley’s Elliptic Functions in India [8], [9, pp. 9–
15]. Narayana Aiyar’s first published problem gives strong evidence that Ramanu-
jan was acquainted with the book. Since Narayana Aiyar and Ramanujan were close
friends, who worked on mathematics together during the evenings, it is inconceivable
that Narayana Aiyar had not told Ramanujan about this book, if indeed Ramanujan
had not already been aware of it before coming to Madras in 1910. Moreover, we also
know that Ramanujan had made most of his voluminous contributions to elliptic func-
tions, in particular to modular equations, in the two or three years prior to his departure
for Cambridge, i.e., during the period when Ramanujan and Narayana Aiyar worked
together.
We commented above on the interest of S. Narayana Aiyar’s problems in geome-
try, and so we offer here two of them. Paul Yiu, in an article under preparation, plans
to discuss all of Narayana Aiyar’s geometrical problems. The first appeared in The
Educational Times on June 1, 1910, but although four solutions were eventually pub-
lished, they did not appear until 1912, indicating that it took a while for readers to
work out solutions. The second appeared on August 1, 1910, and no solution was ever
published.
#16889 If a conic inscribed in a triangle passes through two points which are isogonal
conjugates with respect to the triangle, then the tangents to this conic at these points
intersect either in the incenter or in one of the excenters of the triangle.
#16923 If the polar of the incenter (or an excenter) with respect to an inconic of a
triangle passes through the circumcenter, then the pedal circle with respect to the tri-
angle of the points in which the polar meets the conic touches the nine-point circle of
the triangle.
As with Hardy, several of the initial problems that Narayana Aiyar submitted to The
Educational Times asked for the evaluation of intriguing integrals. In his first problem
(#16238) on integrals, published on July 1, 1907, readers were asked to evaluate
Z π
emr cos θ

−1 r sin θ
1
cos mr sin θ ± n tan dθ.
0 2 2 n a − r cos θ
(a − 2ar cos θ + r ) 2

Here, r < a, presumably n is a positive integer, and no conditions were prescribed for
the parameter m. If one takes the plus sign above, Narayana Aiyar claims that the an-
swer is simply πa −n , but if the minus sign is taken, the evaluation is more complicated.
There were no published solutions. Integrating a suitable function over the unit circle
and using the residue theorem, M. D. Hirschhorn supplied the author with an elegant
solution.
Trigonometric functions were featured in most of Narayana Aiyar’s evaluations of
integrals. We give one further example (slightly reformulated) (#16659), which is more
difficult than the one above. Show that
Z ∞
cos(φ1 + φ2 + · · · + φn ) − θ sin(φ1 + φ2 + · · · + φn ) dθ 1
= π,
0 ρ1 ρ2 · · · ρn 1 + θ2 2
where, for 1 ≤ j ≤ n,
jr sin θ p
φ j = tan−1 and ρj = 1 − 2 jr cos θ + j 2r 2 .
1 − jr cos θ
No solution was ever published, but Hirschhorn has solved this problem and gen-
eralized it; his generalized version appears as problem 11600 in this issue of the
M ONTHLY.
Lastly, we relate the elementary problem (#16187) on probability posed by V. Ra-
maswami Aiyar and solved by S. Narayana Aiyar. Tejashri and Atul play a series of
games in which each game is won by Tejashri or Atul with equal probability, for a fixed
contributed stake on each game. Atul has enough money to play an indefinite number
of games, while Tejashri has no money but pretends that she also has sufficient funds
to play indefinitely. Show that the probability that Tejashri will be able to maintain her
pretension for 2n − 1 games or 2n games is

135 2n − 1 1 2n
··· = 2n .
246 2n 2 n
We now briefly turn to Narayana Aiyar’s questions and solutions in the Journal
of the Indian Mathematical Society. As previously indicated, the Indian Mathemat-
ical Society was originally called the Indian Mathematical Club, and the first three
issues of the organization’s journal, published in the years 1909–1911, were entitled
the Journal of the Indian Mathematical Club. Prior to the publication of the Jour-
nal of the Indian Mathematical Club, problems were evidently published in Progress
Reports. In the sixth, seventh, and eighth reports, Narayana Aiyar solved 4, 6, and 1
problems, respectively. In the first volume of the Journal of the Indian Mathematical
Club, Narayana Aiyar published 11 problems and 15 solutions, with two of the solu-
tions being those for his own problems. It is curious that in future issues there appear
no further solutions by Narayana Aiyar to other submitted problems. Printed with the
first issue are two Progress Reports, in which Narayana Aiyar proposed three addi-
tional problems, while in previous Progress Reports, he had proposed another three
questions. The first proposed problem in the Journal of the Indian Mathematical Club
is Question 59, indicating that the problems section was a continuation of those from
the Progress Reports. Narayana Aiyar was one of seven mathematicians publishing a
large number of problems and solutions in the maiden volume, with the others being
S. Narayanan, N. B. Pendse, V. Ramaswami Aiyar, P. V. Seshu Aiyar, V. R. Sonti,
and K. J. Sanjana. Recall that V. Ramaswami Aiyar founded the Indian Mathematical
Club in 1907 and that P. V. Seshu Aiyar was Ramanujan’s mathematics teacher at the
Government College of Kumbakonam during the latter’s one-year attendance at the

College. Clearly, all seven of these mathematical enthusiasts over a period of several
years had built a storehouse of problems waiting to be published for the enjoyment and
challenge of their contemporaries.
Almost all of the questions submitted by Narayana Aiyar in the initial issue of the
Journal of the Indian Mathematical Club were, not surprisingly in view of our previous
discussion, in geometry. In the next three years, he did not submit any new problems.
He resumed the publication of problems in Volume 5 (published in 1913) of the Jour-
nal of the Indian Mathematical Society. Now his interest turned to series and integrals,
which is not surprising, for he certainly came under the influence of Ramanujan, who
became a clerk in S. Narayana’s office in March, 1912. Ramanujan did not have a keen
interest in geometry, as is clearly demonstrated in both his published papers [13] and
notebooks [14]. However, Ramanujan loved integrals and series! Many of Narayana
Aiyar’s questions in 1915 and 1916 were on series of hypergeometric type, possibly
reflecting Ramanujan’s long-abiding interest in the subject, with Chapters 10 and 11
in his second notebook [14] devoted to the subject. Problems submitted by Narayana
Aiyar in later years were generally in analysis and elliptic functions.
Table 1. Questions by S. Narayana Aiyar.
Volume and Year No. of Proposed Questions
Progress Reports (1908) 6

1(1909) 11
2(1910) 0
3(1911) 0
4(1912) 0
5(1913) 2
6(1914) 4
7(1915) 7
8(1916) 2
9(1917) 1
10(1918) 1
15(1923–24) 3
16(1925–26) 2
17(1927–28) 2
18(1929–30) 1
5. EPILOGUE. Except for a few family pictures, only one photograph of S. Na-
rayana Aiyar is extant. Narayana Aiyar was one of the delegates to the meeting of the
Indian Mathematical Society in Bombay in 1919, and so can be seen in this group
photo. This photograph is extremely valuable, because many important people in the
life of Ramanujan are present. These include (numbering from the left) S. Narayana

Delegates to the Indian Mathematical Society meeting in Bombay, 1919. Narayana Aiyar is third from left in
the first row.
Aiyar (3, first row), P. V. Seshu Aiyar (4, first row), V. Ramaswami Aiyar (5, sec-
ond row), R. Ramachandra Rao (10, second row), and S. R. Ranganathan (2, third
row). Ranganathan wrote the first book biography of Ramanujan in English [15]. He
was the librarian at the University of Madras and is internationally recognized today
as one of the founders of modern library science. Also mentioned in this article are
M. T. Naraniengar (7, second row) and S. Narayanan (9, third row). For a recent ac-
count of the history of the Indian Mathematical Society and of mathematics in India
during the past century, see the article by J. W. Dauben and R. Parikh [11].
We are grateful to V. Viswanathan for providing detailed information about his
grandfather, S. Narayana Aiyar. Viswanathan’s mother was Kamakshi, who entered
the world in 1907 and possibly inspired her father to resume intense mathematical ac-
tivity. Narayana Aiyar’s son, N. Subbanarayanan [16], has written about his father and
his father’s relationship with Ramanujan. In particular, he writes that his father and
Ramanujan exchanged slates prior to the latter’s departure for England. In those days
paper was expensive, and so, in particular, school children performed all of their work
on slates. Before departing for England, Ramanujan did almost all of his mathematics
on his slate. The bestowing of his slate to Narayana Aiyar clearly indicates the respect
that Ramanujan had for him. Subbanarayanan tells that after Ramanujan’s death he
often went to the library to fetch books on elliptic functions for his father. Obviously,
his love for the subject was not abated by the demise of Ramanujan.
Jim Tattersall supplied copies of all the problems and solutions that Narayana Aiyar
submitted to The Educational Times, and also notified the author about the Progress
Reports. Moreover, he informed us of [17], from which we took much of our infor-
mation about The Educational Times. We are enormously beholden to Paul Yiu for
providing solutions to several of Narayana Aiyar’s problems in geometry, and to Mike
Hirschhorn for solutions to some of Narayana Aiyar’s difficult problems on integrals.
We are also grateful to Antreas Hatzipolakis, Clark Kimberling, and John Wetzel for
their comments. The referees also offered several helpful suggestions.

REFERENCES
1. S. Narayana Aiyar, Note on expansion in series, Math. Quest. Sols. from the Educational Times 13 (1907)
49–50.
2. , On the integration of sin mx
sin x , J. Indian Math. Club 1 (1909) 9–10.
3. , The distribution of primes, J. Indian Math. Soc. 5 (1913) 60–61.
4. , Some theorems in summation, J. Indian Math. Soc. 5 (1913) 183–186.
5. , Some infinite determinants, J. Indian Math. Soc. 7 (1915) 51–55.
6. , A small theorem, J. Indian Math. Soc. 8 (1916) 181.
7. B. C. Berndt and R. A. Rankin, Ramanujan: Letters and Commentary, American Mathematical Society,
Providence, RI, 1995; London Mathematical Society, London, 1995.
8. , The books studied by Ramanujan in India, Amer. Math. Monthly 107 (2000) 595–601. http:
//dx.doi.org/10.2307/2589114
9. , Ramanujan: Essays and Surveys, American Mathematical Society, Providence, RI, 2001; Lon-
don Mathematical Society, London, 2001.
10. A. Cayley, An Elementary Treatise on Elliptic Functions, 2nd ed., 1895; reprinted by Dover, New York,
1961.
11. J. W. Dauben and R. Parikh, Beginnings of modern mathematics in India, Current Science 99 (2010)
300–322.
12. S. Ramanujan, On Question 330 of Professor Sanjana, J. Indian Math. Soc. 4 (1912) 59–61.
13. , Collected Papers, Cambridge University Press, Cambridge, 1927; reprinted by Chelsea, New
York, 1962; reprinted by the American Mathematical Society, Providence, RI, 2000.
14. , Notebooks (2 volumes), Tata Institute of Fundamental Research, Bombay, 1957.
15. S. R. Ranganathan, Ramanujan: The Man and the Mathematician, Asia Publishing, Bombay, 1967.
16. S. Subbanarayanan, My father and Ramanujan, in Ramanujan: Letters and Reminiscences, Vol. 1, P. K.
Srinivasan, ed., Muthialpet High School, Madras, 1968, 112–115.
17. J. J. Tattersall, S. L. McMurran, and F. Coughlin, Women and the ‘Educational Times’, Proc. Canad.
Soc. Hist. Philos. Math. 16 (2003) 250–260.
BRUCE C. BERNDT received an A.B. from Albion College, Albion, Michigan, in 1961, and a Ph.D. from
the University of Wisconsin, Madison, in 1966. He joined the faculty at the University of Illinois in 1967 after
one year at the University of Glasgow. Twenty-eight students have received their doctorates under his direction,
with six further students currently writing dissertations under him. Since 1977, he has devoted almost all of his
research efforts toward proving the claims left behind by Ramanujan in his earlier notebooks and lost notebook.
With George Andrews, he is currently preparing a projected five volumes on Ramanujan’s lost notebook, two
of which have been published to date.
Department of Mathematics, University of Illinois, 1409 West Green Street, Urbana, IL 61801, USA
berndt@illinois.edu

On Forming Committees
P. Haxell
Abstract. This paper surveys a collection of results on finding special sets of vertices in graphs
with vertex partitions, all of which can viewed as models for “committee-choosing” problems.
1. INTRODUCTION. As usual, the dean of your institution wishes to form a com-

mittee, consisting of one representative from each department in the faculty. Choosing
such a committee without any restrictions would be an easy task; however, as every-
body knows, there are certain pairs of faculty members who really should not be on a
committee together. Neither one of them can tolerate the other having the last word.
Thus in order to ensure finite committee meetings, the dean should avoid choosing
both members of such a pair when selecting the members of the committee. In this
article we are interested in finding out when this is possible.
First we describe the problem in terms of graphs. Let G be a graph, and suppose the
vertex set of G (the set of faculty members) is partitioned into classes (departments).
Each edge of G is a pair of vertices representing a conflict, i.e., a pair of faculty mem-
bers who should not both be chosen for the committee. Note that since we will choose
only one representative from each department, we may assume that no two vertices
in the same class are joined by an edge. An independent transversal (abbreviated IT)
of G is a set T of vertices in G containing exactly one vertex from each class of the
partition (a transversal) that is independent, meaning no two vertices in T are joined
by an edge. In other words, an independent transversal in G is exactly the notion of a
good choice of committee from the faculty.
It only takes a moment to see the first piece of bad news for the dean: that faculties
can exist that simply do not have any good choice of committee. Imagine for example
a faculty in which no member of the Department of Environmental Studies, however
open-minded, could see eye to eye on any issue with any member of the Department
of Mountaintop-Removal Mining Development. (This corresponds to a complete bi-
partite graph between two partition classes in G.) Thus we might ask whether there is
at least an easy way to decide whether a given partitioned graph has an IT or not.
This brings us to the second piece of bad news. Consider the restricted version of
the committee-choosing problem in which each faculty member is captivated by one
very important two-sided issue. Two people holding opposite views on the same issue
simply can’t be on a committee together. Let us suppose further that no department
contains two members who care about the same issue. This situation then corresponds
to the graph G being a disjoint union of complete bipartite graphs, in which each
complete bipartite graph contains at most one vertex of each partition class. Can we
efficiently determine if a good committee exists? Unfortunately almost certainly not,
as this is a model of the most basic NP-complete problem, the SAT problem (see,
e.g., [15]). Suppose A is a boolean formula in conjunctive normal form, that is, A is
a conjunction of clauses, each of which is a disjunction of propositional variables and
negations of variables. Let G be the graph in which each clause of A corresponds to a
partition class of G, and the vertices of the partition class are labelled with the variables
and negations of variables that occur in the clause (see Figure 1). For each variable x,
November 2011] ON FORMING COMMITTEES 777

join each vertex labelled x to each vertex labelled with the negation of x. (Thus each
“issue” is represented by a variable.) Then in the resulting partitioned graph G, an
independent transversal corresponds exactly to a satisfying truth assignment for A.
ū y z̄ w̄ x ū
u
x̄
y ȳ
z̄
w
ȳ w z
Figure 1. The graph corresponding to the boolean formula A = (u ∨ y ∨ w) ∧ (ū ∨ y ∨ z̄) ∧ (w̄ ∨ x ∨ ū) ∧
(x̄ ∨ ȳ ∨ z̄) ∧ ( ȳ ∨ w ∨ z). The square vertices form an IT corresponding to the truth assignment u = x = T ,
y = z = F, and w can be either T or F.
What help can we then offer to the dean? The best we can do is to provide some
sufficient conditions for an IT to exist in a given vertex-partitioned graph that are not
too hard to check. In the next section we show that if the partition classes are big
enough compared to the maximum degree of the graph, then an IT always exists. In
Section 4 we describe a somewhat more complicated condition for the existence of an
IT, which nevertheless can be applied in certain circumstances. Here we will appeal to
Sperner’s lemma, a result from combinatorial topology.
A very optimistic dean might even want to form many disjoint committees, perhaps
even imagining a perfect world in which the entire faculty could be partitioned into
disjoint committees. This “happy dean” problem is captured by the notion of strong
colouring, which we discuss in Section 3. We end the paper with some remarks and
pointers to other related topics.
2. MAXIMUM DEGREE (LIMITED PERSONAL CONFLICT). In this section

we discover that if no faculty member conflicts with too many others, then the dean can
choose a good committee as long as each department is large enough. Recall that for
a vertex x of a graph G, the degree d(x) of x in G is the number of edges incident to
x. We denote by 1(G) the maximum degree of G; in other words, no faculty member
conflicts with more than 1 others.
Theorem 2.1. Let G be a graph with a vertex partition. Suppose each partition class
has size at least 21(G). Then G has an IT.
Theorem 2.1 first appeared explicitly in [17], and has been applied in a number
of other settings, for example [2, 8, 11, 25]. It is best possible, as Szabó and Tardos
[29] gave constructions of graphs G with partition classes of size 21(G) − 1 that do
not have independent transversals (see also Bollobás, Erdős, and Szemerédi [10], Jin
[23], and Yuster [30] for earlier constructions for certain values of 1). A more precise

version of Theorem 2.1, giving a bound in terms of 1 and the number of partition
classes m, was given in [20] (see also [7]).
We will derive Theorem 2.1 from a more general result, which gives information
about the structure of edge-minimal graphs that do not have independent transversals.
For a graph G with vertex set V (G) and edge set E(G), and a set of vertices W ⊆
V (G), the neighbourhood of W is Ŵ(W ) = {y : wy ∈ E(G) for some w ∈ W }. When
W = {w} we use the abbreviation Ŵ(w) for Ŵ({w}). We say W dominates every vertex
in Ŵ(W ), and we say W dominates G if Ŵ(W ) = V (G). We use the notation G[W ] to
denote the subgraph of G induced by W , that is, the graph with vertex set W and edge
set {x y ∈ E(G) : x, y ∈ W }. For a set Z of edges we write v(Z ) for the set of vertices
incident to edges of Z . Finally for a vertex x in a vertex-partitioned graph we denote
by V (x) the vertex class containing x. Our proof below follows that of [9].
Theorem 2.2. Let G be a graph, and suppose V1 ∪ · · · ∪ Vm is a partition of V (G)

into m independent vertex classes. Suppose G has no IT, but for every edge e the graph
G − e, formed by removing e from G, has an IT. Let e = x y ∈SE(G). Then there exists
a subset S ⊆ {V1 , . . . , Vm } and a set of edges Z of G S = G[ Vi ∈S Vi ] such that
1. V (x), V (y) ∈ S and e ∈ Z ,
2. v(Z ) dominates G S ,
3. |Z | ≤ |S| − 1.
To see that Theorem 2.2 implies Theorem 2.1, let G be as in Theorem 2.1 and
suppose on the contrary that it has no IT. Remove edges one by one from G until the
resulting graph satisfies the assumptions of Theorem 2.2, and let S be the subset of
classes given by Theorem 2.2. The number of vertices that can be dominated by v(Z ),
a set of size at most 2|S| − 2, is at most (2|S| − 2)1, but G S contains 21|S| vertices.
Thus conclusions (2) and (3) cannot hold, so this contradiction shows that G must have
had an IT.
Proof. We prove Theorem 2.2 by induction on m. Let G and e = x y be as in the

statement of the theorem. The assertion of the theorem is trivially true when m = 1,
so assume m ≥ 2 and that the statement is true for smaller values of m.
Choose an IT T of G − e. Then x, y ∈ T , since otherwise T would be an IT of G.
We form a new graph H by
• removing the vertex set W = Ŵ({x, y}) from G (note x, y ∈ W ), and
• unifying the remaining vertices in V (x) ∪ V (y) into one new vertex class Y ∗ (and
removing any edges inside Y ∗ ).
Each class Vi other than V (x) and V (y) just becomes Yi = Vi \ W in H . Note that
each class apart from possibly Y ∗ is nonempty because it still contains an element of T .
Case 1: Y ∗ = ∅.
In this case set S = {V (x), V (y)} and Z = {e}. Then v(Z ) = {x, y} dominates all of
G S as required.
Case 2: Y ∗ 6 = ∅.
First we verify that H does not have an IT. Suppose on the contrary that T ′ is an IT
for H . Let z be the vertex of T ′ in Y ∗ . Then by definition of Y ∗ , in G we have either
z ∈ V (x) or z ∈ V (y). But then in the first case, by definition of H the set {y} ∪ T ′
is an IT of G, and in the second case {x} ∪ T ′ is an IT of G. This contradiction shows
that H has no IT.

Now we remove edges one by one from H until we obtain a graph H ′ with no IT
but such that the removal of any edge results in a graph with an IT. Note that each
vertex w in Y ∗ has degree at least one in H ′ , since otherwise T \ {x, y} ∪ {w} would
be an IT of H ′ . Let e′ be any edge of H ′ incident to a vertex of Y ∗ . Then since H ′ has
m − 1 vertex classes, by induction there exists a set S ′ of vertex classes (containing
Y ∗ ) together with a set of edges Z ′ in HS′ ′ that satisfies the conclusions (1)–(3). Note
that since H ′ was obtained from H by removing edges we also know that v(Z ′ ) dom-
inates HS ′ . But then setting S = S ′ \ {Y ∗ } ∪ {V (x), V (y)} and Z = Z ′ ∪ {e} gives the
required conclusion.
In fact it follows from the above proof that the set Z of edges is a matching, that is,
no two edges in Z share a vertex.
3. STRONG COLOURING (THE HAPPY DEAN PROBLEM). Suppose the

dean does indeed have a good choice of committee. What could be better than finding
many committees, say even a partition of the entire faculty into disjoint committees
(assuming, unrealistically here, that the dean wouldn’t expect any faculty member
to be on more than one committee at once)? This delightful situation for the dean
is related to the notion of strong colouring in graphs. A partition of the vertex set
V (G) into independent sets (committees) I1 , . . . , Ir is called a proper colouring of the
graph G with the r colours 1, . . . , r , a very well-studied notion in graph theory. Here
however the dean also requires something extra, namely that each class (department)
contains exactly one member of each colour (committee); see Figure 2.
Figure 2. A partition into 3 committees (the white, grey, and black committees) for a faculty with 3 depart-
ments.
Let r and n be positive integers such that r divides n, and let G be a graph with n
vertices. We call G strongly r -colourable if for every partition of the vertex set V (G)
into parts Vk of size r , there exists a proper colouring of G with r colours with the
additional property that each Vk contains exactly one vertex of each colour. If r does
not divide n then we say G is strongly r -colourable if the graph obtained by adding
r ⌈n/r ⌉ − n isolated vertices to G is strongly r -colourable. It can be shown that if G is
strongly r -colourable then it is also strongly (r + 1)-colourable. The strong chromatic
number sχ(G) of a graph G is defined to be the minimum r such that G is strongly
r -colourable. This notion was introduced independently by Alon [6] and Fellows [13].
Thus the dean could be guaranteed a partition of the entire faculty G into com-
mittees provided the departments have size at least sχ(G) (and, of course, they all
have the same size). If the unlucky dean finds that the faculty is exactly the graph
given in [29] (see Section 2) with departments of size 21(G) − 1 but no independent
transversal, then not even one committee can be found, let alone a partition into com-
mittees. Thus for some graphs sχ(G) ≥ 21(G). It is conjectured that 21(G) is in
fact also an upper bound for every G, but the best known bounds for strong colouring

in terms of 1(G) are sχ(G) ≤ 31(G) − 1 for every graph [18], and asymptotically
sχ(G) ≤ (1 + o(1))111(G)/4 [19].
Here we give a short proof (by Aharoni, Berger, and Ziv [2], based on [18]) that
sχ(G) ≤ 31(G) for every G. We begin with a slight extension of Theorem 2.1 that
says that in fact, if all departments have size at least 21(G), then the dean can even
choose one specific person on the faculty in advance (who is especially suitable to
chair this committee, say), and still find a good committee that contains this person.
Theorem 3.1. Let G be a vertex-partitioned graph of maximum degree 1 in which

each partition class has size at least 21. Let x ∈ V (G). Then G has an IT contain-
ing x.
Proof. Let V1 , . . . , Vm denote the vertex classes of G (again we may remove any edge
that is inside a class). Assume without loss of generality that x ∈ V1 . Let G ′ be the
graph, with vertex classes Vi′ ⊆ Vi for 2 ≤ i ≤ m, obtained by removing V1 and all
neighbours of x from G. Note that each Vi′ is nonempty since |Vi | ≥ 21. If G ′ has an
IT T then T ∪ {x} is an IT of G as required. If G ′ has no IT then remove edges from G ′
until every remaining edge prevents an IT. Let e be an arbitrary edge of the resulting
graph H (which must exist since each class of H is nonempty). We apply Theorem 2.2
to H and e to obtain a subset S of classes and a set of edges Z of size at most |S| − 1
such that v(Z ) dominates HS . But v(Z ) can dominate at most 21|Z | ≤ 21(|S| − 1)
vertices, and HS contains at least 21|S| − 1 vertices. This contradiction shows that
H has an IT, and hence G has an IT containing x.
We represent a proper colouring of G with colours 1, . . . , r by a function α :

V (G) → {1, . . . , r }, where {x : α(x) = j} is independent for each j.
Theorem 3.2. Let G be a graph of maximum degree 1. Then sχ(G) ≤ 31.
Proof. By adding isolated vertices as necessary, we may assume n = |V (G)| is divisi-

ble by 31. Let a partition P = V1 ∪ · · · ∪ Vm of V (G) into classes of size 31 be fixed,
and suppose on the contrary that there is no suitable colouring of G with respect to
P . Fix a maximum partial colouring (MPC) α of G, that is, a proper colouring using
31 colours of as many vertices of G as possible such that no two vertices in the same
partition class have the same colour. Suppose V1 is a class that contains an uncoloured
vertex x. Then some colour in {1, . . . , 31} is not used in V1 ; let us call this colour
red. For each i, 2 ≤ i ≤ m, let ri denote the red vertex in Vi (if it exists). Our plan
is to “swap colours” between each ri and another vertex in the same class as ri . For
2 ≤ i ≤ m set Wi = Vi \ {v : α(v) = α(z) for some z ∈ Ŵ(ri )} (if ri does not exist
then Wi = Vi ), so Wi is the set of vertices whose colour could be given to ri . Set also
W1 = V1 . Then each |Wi | ≥ 21, so by Theorem 3.1 the graph G[W1 ∪ · · · ∪ Wm ] has
an IT T with x ∈ T . Modify α by giving colour red to every vertex of T , and for each
i for which ri exists and ri ∈ / T , give ri the colour of the element ti of T in Vi . Then
since T was an IT, and by definition of the Wi , this gives a valid colouring α ′ . More-
over α ′ colours all the vertices that were coloured by α together with x, contradicting
the fact that α was an MPC. Therefore a valid colouring must exist for P , and thus we
conclude sχ(G) ≤ 31.
Finding the correct function of 1(G) that bounds sχ(G) from above seems to
be difficult (see Problem 4.14 in Jensen and Toft [22]). For example, it is not even
known whether sχ(G) ≤ 4 for every graph of maximum degree two. By [18] we know
sχ(G) ≤ 5 for such graphs.

4. A HALL-TYPE CONDITION. The criterion for the existence of an IT that we
describe in this section uses a rather more complicated notion than maximum de-
gree or class size, so we begin with the relevant definition. Let G be a graph, and
suppose V1 ∪ · · · ∪ Vm is a partition of V (G) into m independent vertex classes. Let
. . . , Vm } be a subset of classes. We call an independent set I S of vertices of
S ⊆ {V1 , S
G S = G[ Vi ∈S Vi ] special for S if for every independent subset J of vertices of G S
with |J | ≤ |S| − 1, there exists v ∈ I S such that J ∪ {v} is also independent. We can
think of a special independent set as a “neutral team” for the set S of departments: a
group of mutually nonconflicting members from those departments that can augment
any “small” set of mutually nonconflicting members. (Here “augment” is not quite the
right word, since if I S actually contains an element v of J then J ∪ {v} is of course
the same size as J and satisfies the condition.) Note that the departments from which
the members of I S or J are taken do not figure in this definition.
The reason we refer to the upcoming theorem as a Hall-type condition for the ex-
istence of an IT is by analogy with Hall’s classical theorem for the existence of a
matching of size |A| in a bipartite graph H with vertex classes A and X (here bipartite
means V (H ) = A ∪ X , and A and X are disjoint independent sets in H ).
Theorem 4.1. (Hall’s Theorem) A bipartite graph H with vertex classes A and X has
a matching of size |A| if and only if for every subset S of A we have |Ŵ(S)| ≥ |S|.
We may now state the main theorem of this section, which is from [4]. Note that
in contrast to Hall’s theorem this is only a one-way implication; for an if-and-only-if
version see Section 5 (where we also show that this theorem implies Hall’s theorem).
Theorem 4.2. Let G be a graph with a vertex partition into independent classes. Sup-
pose that for every subset S of classes, the graph G S contains an independent set I S
that is special for S. Then G has an IT.
The dean may be justifiably dubious that this theorem could possibly be useful,
as checking whether a partitioned graph satisfies the condition (of having a special
independent set for every subset of classes) looks hopelessly complicated. Thus to
convince the dean to read on, we first give a quick application of Theorem 4.2 to a well-
known combinatorial problem. Then in the following subsection we will describe the
proof, which uses Sperner’s lemma. We mention that another application of Theorem
4.2 (and in fact its original motivation) is an extension of Hall’s theorem to hypergraphs
(see [4]).
4.1. Application of Theorem 4.2. Our application is related to the famous “cycle
plus triangles” problem, popularised by Erdős in the 1980s, and finally solved by Fleis-
chner and Stiebitz [14] in 1992 and with a different proof by Sachs [28] in 1993. It asks
whether every union of a cycle C3k of length divisible by three, together with a set of k
disjoint triangles on the same vertex set, has a proper colouring with three colours (see
Figure 3). In the language of Section 3, this is the same as asking whether every C3k
is strongly 3-colourable. This question of Erdős was motivated by an earlier question
of Du, Hsu, and Hwang [12], who asked whether every C3k with an arbitrary vertex
partition into classes of size three has an IT. This problem was also unsolved until the
proofs of Fleischner and Stiebitz, and Sachs, of the stronger statement. Both of these
proofs are ingenious but quite difficult. Here we can give a solution to the question of
Du, Hsu, and Hwang that is almost immediate from Theorem 4.2, and in fact is more
general.

Figure 3. A 3-coloured cycle-plus-triangles graph on the cycle C9 . Compare to Figure 2, which is a different
drawing of C9 , with the same three triangles represented as departments.
Theorem 4.3. Let G be a graph of maximum degree two, in which each cycle is of
length divisible by three. Then G has an IT with respect to any vertex partition into
classes of size at least three.
Proof. Fix a vertex partition, and remove any edge that is inside a class. Let S be any
subset of classes. Then the graph G S has at least 3|S| vertices, and its components are
cycles of length divisible by three and paths. Let I S be a maximum independent set in
G S that is spaced three apart. For example, I S could be formed by taking, from each
path component, an end vertex and then every third vertex starting from that end, and,
from each cycle component C3k , an independent set of size k containing every third
vertex. We claim that I S is special for S. To see this, let J be an arbitrary independent
set in G S of size |S| − 1. Since I S is spaced three apart, each vertex of J is adjacent
to at most one vertex of I S . Thus since |I S | ≥ |S|, there is a vertex v ∈ I S such that
J ∪ {v} is independent. Thus I S is special, and so by Theorem 4.2 we know that G has
an IT.
4.2. Proof of Theorem 4.2. First we recall Sperner’s lemma, a basic result from com-
binatorial topology. Suppose F is a triangulation of the (m − 1)-dimensional simplex
6m−1 . For each point x of F we denote by f (x) the face of 6m−1 containing x in its
interior. A labelling of the points of F with elements of {1, . . . , m} is called a Sperner
labelling if
• each vertex of 6m−1 receives a different label, and
• each point x of F receives the same label as some vertex of f (x).
An example of a Sperner labelling is shown in Figure 4. We call a simplex of the
triangulation fully-labelled if it receives all m labels on its vertices.
Theorem 4.4. (Sperner’s Lemma) Let F be a triangulation of 6m−1 with a Sperner

labelling. Then the number of fully-labelled simplices in F is odd.
As is usually the case, the implication of Sperner’s lemma that is useful for us is that the
number of fully-labelled simplices is nonzero. Note that there are three fully-labelled
simplices in Figure 4.
For the proof of Theorem 4.2 we also need to know that certain special triangula-
tions of 6m−1 exist (see Figure 4). The 1-skeleton of a triangulation F is the graph
whose vertices are the points of F, and whose edges are the 1-dimensional simplices
of F.

1
1 2
1 2
1 2
3
2
1
2 1
2
2
3 2
2 3
Figure 4. A Sperner labelling of a triangulation F of 62 with the properties in Lemma 4.5.
Lemma 4.5. There exists a triangulation F of 6m−1 with the following properties.
(i) If x and y are adjacent in the 1-skeleton of F then one of f (x) and f (y)
contains the other,
(ii) if x has neighbours in the 1-skeleton of F on the boundary of f (x), then these
neighbours are the vertices of a simplex of F.
A proof of Lemma 4.5 can be found in [4]. A different construction was given in [3].
a A
b B
c C
d D
V1 α V2
δ
V3
Figure 5. The graph G.
Proof. We may now give the proof of Theorem 4.2. Let the graph G with vertex parti-
tion V1 ∪ · · · ∪ Vm be given (see Figure 5 for an example). Let F be the triangulation of
6m−1 given by Lemma 4.5, and let z 1 , . . . , z m denote the vertices of 6m−1 . We define a
function g that assigns to each point of F a vertex of G with the following properties.

(1) For each point x, the vertex g(x) is an element of I S where S = {Vi : z i is a
vertex of f (x)}.
(2) If points x and y are adjacent in the 1-skeleton of F then g(x) and g(y) are not
joined by an edge in G.
(3) The labelling ℓ of F defined by ℓ(x) = i where g(x) ∈ Vi is a Sperner labelling.
A suitable function g for the graph G in Figure 5 is shown in Figure 6.
b
D C
c
c A
A
d
d
A D
A
A
D
Figure 6. The triangulation F with the vertex g(z) of G assigned to each point z.
Our motivation here is as follows. Once we have found g, by (3), the labelling ℓ
is a Sperner labelling. (The labelling ℓ for the function g in Figure 6 is the labelling
in Figure 4.) Therefore by Theorem 4.4 there exists a simplex W of F that gets all m
labels on its vertices (e.g., the shaded triangle in Figure 6). By (2), the set of vertices
I = {g(x) : x ∈ W } is independent in G. By definition of ℓ, each v ∈ I lies in a distinct
class Vi . Therefore I is an IT of G as required. (I = {c, A, γ } in our example.)
Thus to finish the proof, we define g on the points x of F, in increasing order of
dim( f (x)). For the points in zero-dimensional faces (i.e., the vertices z 1 , . . . , z m of
6m−1 ) we choose an arbitrary vertex vi ∈ I{Vi } (which is nonempty by the assumption
of the theorem) and set g(z i ) = vi . This satisfies (1), satisfies (2) by Lemma 4.5(i), and
is consistent with (3).
Now suppose j ≥ 1 and we have defined g on all points in faces of dimension
smaller than j, and possibly some in faces of dimension j, such that (1)–(3) are satis-
fied (see Figure 7). Let x be in a face f of dimension j (e.g., the grey point in Figure 7),
and let S = {Vi : z i is a vertex of f }. (S = {V1 , V2 , V3 } in Figure 7.) By Lemma 4.5(i)
and (ii), if x has any neighbours in the 1-skeleton of F in faces of dimension smaller
than j, then the set U of such neighbours lies in the boundary of f and forms a simplex
of F. Thus |U | ≤ dim( f ) = |S| − 1 and by (2) J = {g(x) : x ∈ U } is independent in
G. (J = {b, γ } in Figure 7.) Thus J is a “small” independent set, and so the special in-
dependent set I S contains a vertex v such that J ∪ {v} is independent. (In our example
suppose I S = {c, d, A, D}. Then v = c is a suitable choice.) Then we set g(x) = v, so
(1) and (3) are satisfied. To check (2), if y is a neighbour of x in a smaller dimensional
face then (2) is satisfied by choice of v. If y is in a face of dimension j (e.g., the white
point in Figure 7) then by Lemma 4.5(i) it must be in f as well. Thus if g(y) has

a
b
D C
c
A
D
Figure 7. Defining g on the points of F.
already been defined then g(y) ∈ I S by (1) (g(y) = c in our example) and therefore is
not adjacent to g(x) in G. This completes the definition of g, and hence the proof.
5. REMARKS. The if-and-only-if version of Theorem 4.2 is as follows.
Theorem 5.1. Let G be a graph, and suppose V1 ∪ · · · ∪ Vm is a partition of V (G)

into m independent vertex classes. Then G has an IT if and only if the following holds:
forSeach S ⊆ {V1 , . . . , Vm } there exists an independentSset I S of vertices of G S =
G[ Vi ∈S Vi ] such that for every independent subset J ⊆ U ⊂S IU with |J | ≤ |S| − 1,
there exists v ∈ I S such that J ∪ {v} is also independent.
The “only-if” implication is easy: just take I S = I ∩ V (G S ) for an IT I . The proof

of the “if” implication is the same as the proof of Theorem 4.2. To see that Hall’s
theorem is a special case, suppose a bipartite graph H as in Theorem 4.1 is given, such
that |Ŵ(S)| ≥ |S| for every S ⊆ A. Define a graph G whose vertex set is the set of
edges of H by joining e and f by an edge of G if and only if e and f are incident
to the same vertex of X in H . The vertices of G are then partitioned into |A| classes,
according to the vertex of A they are incident to in H . Thus an IT of G is precisely a
matching of size |A| in H . To verify the assumption of Theorem 5.1, let S be a subset
of the vertex classes. Choose I S to be a set of |S| edges of H , all incident to distinct
vertices in X . This is possible by Hall’s condition on H . Then for any set J of edges
of H of size at most |S| − 1, some element e of I S has a different X -vertex from all
elements of J , and thus J ∪ {e} is also independent. Therefore G has an IT.
The use of topological arguments for transversal-type problems has been taken
much further; see, e.g., [26], [29], and [1]. In [1], Aharoni and Berger obtain wide-
ranging results in the much more general setting of matroids, which have many inter-
esting applications.
The dean now knows that good committees exist in many situations, but suddenly
wants to know something else: how do we actually find these committees? Are there
efficient algorithms for finding them? Each of the proofs of the theorems we have
seen (Theorems 2.2, 3.2, and 4.2) does in fact give a procedure for finding the IT that
the theorem claims exists. Unfortunately though, in each case, the number of steps

is potentially exponential in the size of the graph. Thus, sadly, we can’t reassure the
dean that finding these committees is always an easy task. In general, the problem of
finding structures that are guaranteed to exist but are apparently hard to find gives rise
to a fascinating branch of complexity theory; see for example [27].
ACKNOWLEDGMENTS. The author was partially supported by NSERC.
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PENNY HAXELL received her B.Math from the University of Waterloo in 1988 and her Ph.D. from the
University of Cambridge in 1993. She has been a faculty member in the Combinatorics and Optimization
department at the University of Waterloo since 1993. This article is based on an MAA invited address at the
AMS-MAA Joint Meetings in New Orleans in 2007. The author has no good excuse for taking this long to
write it.
Department of Combinatorics and Optimization, University of Waterloo, Waterloo Ont. Canada N2L 3G1.
pehaxell@math.uwaterloo.ca
Mathematics Is . . . and Is Not . . .
“Mathematics is abstract thought, mathematics is pure logic, mathematics is cre-

ative art.”
Paul Halmos, Mathematics as a creative art,
American Scientist 56(4) (1968) 380.
“Mathematics is not a deductive science—that’s a cliché. When you try to prove

a theorem, you don’t just list the hypotheses, and then start to reason. What you
do is trial and error, experimentation, guesswork.”
Paul Halmos, I Want to Be a Mathematician: An Automathography,
Springer-Verlag, New York, 1985, p. 321.
—Submitted by Carl C. Gaither, Killeen, TX

Complete Dissections: Converting Regions
and Their Boundaries
Tom M. Apostol and Mamikon A. Mnatsakanian
Abstract. Classical dissections convert any planar polygonal region onto any other polygonal
region having the same area. If two convex polygonal regions are isoparametric, that is, have
equal areas and equal perimeters, our main result states that there is always a dissection, called
a complete dissection, that converts not only the regions but also their boundaries onto one
another. The proof is constructive and provides a general method for complete dissection using
frames of constant width. This leads to a new object of study: isoparametric polygonal frames,
for which we show that a complete dissection of one convex polygonal frame onto any other
always exists. We also show that every complete dissection can be done without flipping any
of the pieces.
Dedicated to the memory of Martin Gardner
1. INTRODUCTION. DISSECTIONS INVOLVING BOUNDARIES. It is well

known that any planar polygonal region can be dissected into smaller polygonal pieces
that can be rearranged to form any other polygonal region of equal area (see for ex-
ample [4, p. 221]).
What happens to the boundaries in these standard dissections?
Figure 1a shows a dissection of a triangle onto a rectangle, and Figure 1b shows
a dissection of one rectangle onto another rectangle of prescribed altitude. In both
examples, part of the boundary of one polygon ends up inside the other. This is to
be expected, because the initial and final shapes have different perimeters. Figure 20a
shows an even more dramatic example, a famous hinged dissection of Dudeney in
which the entire boundary of a square ends up inside a triangle.
We are interested in dissections that not only preserve areas but also convert the
boundaries onto each other as well. We call these complete dissections. They require
that both regions have equal areas and equal perimeters. Such regions, called isopara-
metric, were introduced in [1], and provide a natural motivation for this paper. To the
best of our knowledge, such dissections have not been previously treated.
(a) (b)
Figure 1. Dissection (a) of a triangle onto a rectangle; (b) of a rectangle onto another rectangle of prescribed
altitude. Dark lines show how the boundaries are transformed.
In Figure 2, the triangle and rectangle are isoparametric, but the dissection shown
is not complete because one boundary is not converted entirely onto the other. In fact,
there is no reason to expect that a complete dissection exists even though the figures are
isoparametric. Nevertheless, our first theorem reveals a surprising and profound result:
for any two isoparametric convex polygonal regions, a complete dissection always
exists. Moreover, the proof shows how to construct such a dissection.
November 2011] COMPLETE DISSECTIONS 789

2
5
2 2
6 6
Figure 2. Standard dissection converting a triangle onto an isoparametric rectangle.
Before proceeding further, the reader might try to find a complete dissection that
converts the triangle in Figure 3 onto the isoparametric isosceles trapezoid shown.
This is a simpler task than for the pair in Figure 2.
Figure 3. Equilateral triangle and an isoparametric trapezoid.
Figure 4 shows a complete dissection of any isosceles triangle onto an isosceles

trapezoid, and of a general triangle onto a trapezoid with the same base angles.
(a)
angle bisector
a A A
C A
B b B
(b) a b
B b C a C
Figure 4. Complete dissection converting (a) isosceles triangle onto isosceles trapezoid, and (b) any triangle
onto a trapezoid. In (a), the unshaded piece has been flipped. In (b) one piece is flipped and divided into two
right triangles, the smaller of which is flipped again.
2. COMPLETE DISSECTION OF POLYGONAL REGIONS. We turn now to

the first principal result of this paper.
Theorem 1. Any two isoparametric convex polygonal regions can be converted onto
one another by complete dissection.
Proof. Consider two isoparametric convex polygonal regions A and B. Figure 5 shows
an example with A triangular and B quadrilateral. This example displays all the essen-
tial features required in treating general convex isoparametric polygonal regions. The
method of proof is suggested by an oversimplified intuitive idea: Remove each bound-
ary and perform a standard dissection of the interior of A to produce the interior of B.
Then restore the two boundaries to obtain the complete dissection.

A B
Figure 5. Isoparametric triangular and quadrilateral regions.
To make this intuitive idea rigorous, refer to Figure 6a, which shows a “frame” of
some constant width w protruding into each region along its boundary. Choose w small
enough so that the inner boundary of each frame will be a simple closed polygon with
the same number of sides. Each region now consists of two parts: the frame plus the
interior region surrounded by the frame. Keep in mind that:
The sum of areas, frame plus interior, is the same for both regions A and B.
Now unfold each frame at the outer vertices (thought of as hinges) and lay it out
horizontally, as shown in Figure 6b. To be specific, use angle bisectors to cut each
frame into trapezoidal pieces with isosceles triangular gaps between adjacent pieces.
(a)
w w
(b) w
Figure 6. Frames unfolded to form adjacent trapezoidal pieces with triangular gaps.
Next, use standard dissections (as in Figure 1) of the regions interior to the frames
in Figure 6a to convert them onto two rectangles with common altitude w (the frame
width), as indicated in Figure 7a. Each rectangle has area equal to that of the interior
that produced it, but, of course, the two rectangular areas are not equal to each other.
(a)
(b)
(c)
Figure 7. (a) Dissection of interior regions produces two rectangles of equal altitude, but not of equal area.
Removing shaded triangular pieces in (b) to fill the gaps as shown in (c) leaves two unshaded congruent
rectangles in (b).
From each of these rectangles, remove smaller triangles as needed (see Figure 7b)
to fill the triangular gaps in Figure 6b. This transforms the unfolded frames into two
congruent rectangles shown in Figure 7c, whose lower bases are the unfolded bound-
aries of A and B. By overlapping these rectangles we obtain their common dissection
(dissection 1) which also converts the full boundary of A onto the full boundary of B.

The leftover unshaded rectangles in Figure 7b are also congruent because the sum of
areas, frame plus interior, is the same for both regions in Figure 6a. By overlapping
the unshaded rectangles in Figure 7b, we obtain a common dissection (dissection 2)
of the two interior dissections inherited from Figure 7a. The union of common dissec-
tions 1 and 2 gives a complete dissection of region A onto region B and completes the
proof.
Obviously, these dissections can be done in many different ways. For example, we
can choose a different frame width w.
3. COMPLETE DISSECTION OF POLYGONAL FRAMES. The frames of con-

stant width used in the proof of Theorem 1 lead us to consider the problem of complete
dissection of such frames, where now both inner and outer boundaries are subject to
conversion. Although frames are more complicated objects than those in Theorem 1,
our second principal result (Theorem 2) states that any two isoparametric frames can
also be converted onto one another by complete dissection. Before discussing Theorem
2, we explain precisely what we mean by a polygonal frame.
In this paper, the term polygonal frame refers to a frame of constant width. It has
parallel inner and outer boundaries with constant distance separating the parallel edges.
We restrict our discussion to convex polygonal frames, that is, frames in which both
the inner and outer polygons are convex.
More precisely, start with a convex n-gon as inner boundary, and any frame width
w. Draw lines outside the n-gon parallel to its sides at distance w from the sides.
Segments of the lines will form another convex n-gon outside the inner boundary, like
the examples in Figure 6a. The frame consists of the region between the two n-gons,
including both boundaries. When unfolded at the outer vertices, the frame forms a set
of adjacent trapezoidal regions akin to those in Figure 6b.
Complete dissections require isoparametric frames, which have equal areas and
equal total perimeters (inner plus outer). Figure 8 shows a charming example (taken
from [1]): a Pythagorean 3:4:5 triangular frame of constant width w = 2, and a square
frame of the same constant width.
15
9 5
3
w= 2
4
12
8 4
w= 2
common dissection
Figure 8. Complete dissections of two special isoparametric frames.

Note that both the areas and total perimeters of these frames have the same numer-
ical value, 48. This is not merely a coincidence, but is a consequence of the following
lemma when w = 2.
Lemma 1. For any convex polygonal frame, its width w, area A, and total perimeter
P are related by the equation
A = 12 Pw. (1)
The proof follows easily by applying the area formula for adjacent trapezoids form-
ing the frame.
As an immediate consequence of Lemma 1 we have the following crucial result:
Corollary 1. All isoparametric convex polygonal frames have the same width.
Now return to Figure 8, which shows an unfolding of each frame into trapezoids of
constant altitude w = 2, followed by a dissection onto a rectangle of the same altitude
with two horizontal bases, the sum of whose lengths is the total perimeter of the frame.
Because the two frames are isoparametric, so are the two rectangles. At the bottom of
Figure 8, the dissected rectangles are superimposed to obtain a common dissection of
the two frames that converts their total boundaries, indicated by the heavy lines.
The same type of argument can be applied to any pair of isoparametric convex
polygonal frames to prove the following theorem.
Theorem 2. Any two convex isoparametric polygonal frames can be converted onto
one another by complete dissection.
4. COMPLETE DISSECTIONS WITHOUT FLIPPING. In the foregoing dissec-

tions, some pieces may have been flipped. Although flipping might reduce the number
of dissection pieces, in some applications, such as skin grafting or laying out carpet-
ing, flipping is undesirable. It is known that every standard dissection can be carried
out without flipping. Now we state our third principal result:
Theorem 3. Every complete polygonal dissection can be done without flipping.
Flipping a piece turns it into its mirror image, so the proof of Theorem 3 reduces
to showing that a complete dissection of a general polygonal piece onto its mirror
image can be done without flipping and in such a way that each edge of the polygon
is converted to the corresponding edge of the mirror image. Such a dissection requires
more than completeness, and we call it strongly complete. Now we will prove that
every polygon can be converted onto its mirror image by using strongly complete
dissection.
First dissect the polygonal piece into triangles, as illustrated in Figure 9. Then per-
form a strongly complete dissection on each triangle.
Figure 9. Any polygonal piece can be cut into triangles.

(a) (b)
Figure 10. (a) Cutting a general triangle into right triangles, each of which has only one leg subject to conver-
sion. (b) Converting one leg of a right triangle onto its mirror image.
To do this, we can dissect each triangle into six right triangles as shown in Figure
10a, where now each right triangle has only one leg as part of the boundary of the
original triangle. This reduces the problem to that of dissecting a right triangle onto its
mirror image, without flipping, so that one leg gets converted onto its mirror image.
Figure 10b shows how a right triangle can be cut into two isosceles triangles. Rotate
one of them to convert the given leg as shown. This completes the proof of Theorem 3.
By combining Theorems 1 and 3 we conclude that any two isoparametric convex
polygons can be converted onto one another by complete dissection without flipping.
Examples. Figure 11 (taken from [1]) shows how easy it is to produce an endless
supply of incongruent isoparametric pairs in which boundaries are transformed onto
one another. In each case, a chord bisects the first region and one piece is flipped, thus
giving a simple but complete dissection of one onto the other. It is reassuring to realize
that every example of such a dissection can also be done without flipping.
isosceles triangle and parallelogram parallelogram and kite regular pentagon and hexagon
Figure 11. Isoparametric regions with boundaries converted onto one another.
Figure 12 shows two different complete dissections converting two isoparametric

curvilinear regions onto one another. In the left figure (taken from [1]), a chord bi-
sects the oval and one piece is flipped as in the examples of Figure 11 to produce the
symmetric heart-shaped figure. A complete dissection without flipping is shown on the
right, where cuts are made along the diagonals of the square inscribed in the oval, and
then the pieces are rotated as shown. This dissection works for any oval having two
perpendicular axes of symmetry.
with flipping without flipping

Figure 12. Two complete dissections of isoparametric curvilinear regions.
Figure 13 shows a complete dissection without flipping of a rectangular frame onto

an isoparametric square frame. A similar dissection also works for any two rectangular
frames of equal width having equal outer perimeters.

A
A
3 1 3 4 2
2
5 B
4 B
Figure 13. Complete dissection of rectangular frame and isoparametric square frame.
5. COMPLETE DISSECTIONS USED TO APPROXIMATE CURVILINEAR

REGIONS. Figure 14 shows a frame that is partially polygonal and partially curvilin-
ear. This example consists of a quadrilateral portion together with a curvilinear portion
obtained as a limit of a portion of a polygonal frame of the same constant width w.
w
w
w w
w w
w
w
Figure 14. Frame with partially curvilinear boundaries, and a circumgon.
The second region in Figure 14 is a circumgon of inradius w that can be regarded

as a degenerate case of the frame on the left, with the inner boundary replaced by a
single point and the curved portion being a circular arc. Circumgons were introduced
in [2], where it was shown that the area A and perimeter P of any circumgonal region
with inradius w are related by the equation A = Pw/2, which is (1) in Lemma 1.
More general frames, partially curvilinear and partially polygonal, can be intro-
duced similarly, so that they share formula (1) of Lemma 1. Then Corollary 1 implies
that any two isoparametric partially curvilinear frames have the same width.
Figure 15 shows an example in which the curvilinear part is semicircular and the
polygonal part is rectangular. This partially semicircular frame and the square frame
shown adjacent to it are isoparametric. Actually, this square frame is only one of in-
finitely many rectangular frames of width w = 2 isoparametric to this semicircular
w
outer edge length
r (r + 1) / 4 + (r + 6) / 2
w=2 w=2
Figure 15. A partially semicircular frame and an isoparametric square frame.

frame. If the inner radius of the semicircular part is r , the common value of both the
area and the perimeter is given by A = P = 2π(r + 1) + 4(r + 2).
Further examples. Figure 16a (borrowed from [1]) shows two isoparametric circular
sectors that are also isoparametric to the same rectangle. Figure 16b shows a circular
frame and an isoparametric square frame of the same width.
2 1
2 1 1
2 3 −1
4
1 3 +1
(a) (b) 4
Figure 16. Partially circular regions isoparametric to polygonal regions.
1 2 1
1 2
Figure 17. Radial slicing of two special circular sectors with their boundaries, rearranged differently to ap-
proximate the same isoparametric rectangle.
Figure 17 shows how the circular sectors in Figure 16a can be dissected into an
even number of radial slices that can be rearranged to form a figure approximating
the same rectangle. In one case the slices are arranged in horizontal layers, and in
the other case in vertical layers. As the number of slices increases without bound,
both approximations have the same rectangle as a limit, with the curved boundaries
becoming the rectangular boundaries, in the style of Archimedes. Figure 18 shows a
dissection of the circular frame in Figure 16b that approximates the square frame. The
Figure 18. Radial slices of a circular frame and its boundary, rearranged to approximate the isoparametric
square frame.

1
1
...
1 1 1
Figure 19. Regular polygonal frames and a circular frame, all isoparametric.
two frames in Figure 16b are members of an infinite family of isoparametric regular
frames of constant width w = 1, examples of which are shown in Figure 19.
Each inner regular n-gon in this family has an incircle of diameter dn , where
3π
dn = − 1,
n tan πn
and each outer regular n-gon has an incircle of diameter Dn = dn + 2w = dn + 2.
6. CONCLUDING REMARKS. This paper introduces a new development in classi-

cal dissection problems: complete dissections that convert regions and their boundaries
onto one another. This requires isoparametric regions, objects of independent interest
introduced in [1].
The methods of this paper can be adapted to more general situations, illustrated
by the example in Figure 20b, a square and equilateral triangle with equal areas and
unequal perimeters. Using a method similar to that for proving Theorem 1, with the
frames in Figure 6 replaced by frames like those in Figure 20c, it is possible to dissect
a general polygonal region onto any other polygonal region of equal area and unequal
perimeter, with the shorter boundary converted onto part of the longer boundary, as in-
dicated by the heavy lines. This suggests further modifications of classical dissections
that might provide interesting applications to recreational mathematics.
(a) (b) (c)

Figure 20. In (a) and (b), the areas are equal and the perimeters are unequal. (c) Modifying the frames in
Figure 6 to treat dissections of the type in (b).
In view of Max Dehn’s counterexample that settled Hilbert’s third problem, exten-
sions of our ideas to dissections of solids in 3-space are not always possible. Never-
theless, circumsolid shells, introduced in [3], have properties analogous to polygonal
frames. For example, all circumsolid shells have constant thickness (see Theorem 8a
in [3]). The analog of Lemma 1 is given by Theorem 10 in [3].
Frederickson’s book [4] gives an admirable introduction to the field of geometric
dissections, and contains a valuable bibliography of known results. Although there is a
vast literature on standard polygonal dissection, we were not able to find any references
relating to complete dissections of the type discussed in this paper.

It is surprising that complete dissections have not been previously discussed in con-
nection with cake slicing. Here’s a natural question: Can we cut a cake with white
icing on top and chocolate icing on its outer edges and rearrange the pieces to form
a cake of another shape (as in Figure 12) so that the white icing stays on top and the
chocolate icing stays on the outer edges of the rearranged cake? Theorem 3 shows that
this can always be done for polygonal cakes.
ACKNOWLEDGMENTS. The authors wish to thank the referees for valuable suggestions that improved
this paper. They also wish to thank Dan Velleman for suggesting a way to considerably shorten our proof of
Theorem 3, and for suggesting the simple example of isoparametric frames in Figure 13.
REFERENCES
1. T. M. Apostol and M. A. Mnatsakanian, Isoperimetric and isoparametric problems, Amer. Math. Monthly
111 (2004) 118–136. http://dx.doi.org/10.2307/4145213
2. , Figures circumscribing circles, Amer. Math. Monthly 111 (2004) 853–863. http://dx.doi.
org/10.2307/4145094
3. , Solids circumscribing spheres, Amer. Math. Monthly 113 (2006) 521–540. http://dx.doi.
org/10.2307/27641977
4. G. N. Frederickson, Dissections: Plane and Fancy, Cambridge University Press, Cambridge, 1997.
TOM M. APOSTOL joined the Caltech faculty in 1950 and became professor emeritus in 1992. He is director
of Project MATHEMATICS! (http://www.projectmathematics.com), an award-winning series of videos
he initiated in 1987. His long career in mathematics is described in the September 1997 issue of The College
Mathematics Journal. He is currently working with colleague Mamikon Mnatsakanian to produce materials
demonstrating Mamikon’s innovative and exciting approach to mathematics.
California Institute of Technology, 253-37 Caltech, Pasadena, CA 91125
apostol@caltech.edu
MAMIKON A. MNATSAKANIAN received a Ph.D. in physics in 1969 from Yerevan University, where he
became professor of astrophysics. As an undergraduate he began developing innovative geometric methods
for solving many calculus problems by a dynamic and visual approach that makes no use of formulas. He is
currently working with Tom Apostol under the auspices of Project MATHEMATICS! to present his methods in
a multimedia format.
California Institute of Technology, 253-37 Caltech, Pasadena, CA 91125
mamikon@caltech.edu

A “Bouquet” of Discontinuous Functions for
Beginners in Mathematical Analysis
Giacomo Drago, Pier Domenico Lamberti, and Paolo Toni
Abstract. We present a selection of a few discontinuous functions and we discuss some peda-
gogical advantages of using such functions in order to illustrate some basic concepts of math-
ematical analysis to beginners.
1. INTRODUCTION. In this paper we present a selection of several discontinuous

real-valued functions of one real variable which we believe could be proposed to any
beginner in mathematical analysis, even to students of secondary school dealing with
the first notions of calculus. Some of them are elementary and well known, others a bit
more sophisticated. Most of them are modeled on the Dirichlet function.
Our aim is to point out some pedagogical advantages of using discontinuous func-
tions rather than classical analytic functions. Usually, young students tend to think
of basic mathematical analysis as a set of rules of calculus and, surely, many under-
graduate students would claim to feel more acquainted with the notion of derivative
rather than with the notion of function. This is not surprising at all, if we consider that
the historical evolution of the notion of function has been very long and troubled;
see Youschkevitch [15] and Kleiner [8, 9]. See also Nicholas [11], Deal [3], and
Thurston [13] for an interesting discussion about some pedagogical issues related to
the definition of function. A modern definition of function can be found, e.g., in the
classic book of Bourbaki [2], published more than two hundred years after the defini-
tion of Johann Bernoulli (1718). As is pointed out in [15, p. 79], a significant step in
this process was the formulation of A. Cournot (1841), which we report here for the
convenience of the reader:
We understand that a quantity may depend on another [quantity], even in case

the nature of this dependence is such that it cannot be expressed in terms of a
combination of algebraic symbols.
This level of generality is commonly attributed to Dirichlet, who in 1829 proposed his
celebrated function D defined on (0, 1) as follows:
(
1, if x ∈ (0, 1) ∩ Q,
D(x) = (1)
0, if x ∈ (0, 1) \ Q.
In fact, this example opened a door to a new world: functions are not just formulas,
or analytic expressions, as was commonly assumed in the 18th century. Functions can
be defined by very general laws.
From a pedagogical point of view, deciding the level of generality of functions to
use with young students is not straightforward. According to Kleiner [9, pp. 187–188],
it is possible to teach an elementary model of analysis by placing emphasis solely on
curves and the equations that represent them, without necessarily talking about func-
tions. Kleiner argues that students would find curves more natural than functions and
November 2011] A “BOUQUET” OF DISCONTINUOUS FUNCTIONS 799

teachers should introduce the notion of function only when there is an evident need for
it. Teaching mathematics should follow the historical evolution of mathematics itself:
new definitions and concepts were introduced when the need arose, and it should be so
in teaching too. Thus, the definition of function should be given gradually, following
its historical path. At first, one should introduce functions as formulas, then as rules,
and only at last, if required, as a set of ordered pairs (after all, as Kleiner [9] points
out, “giving this latter definition and proceeding to discuss only linear and quadratic
functions makes little pedagogical sense”).
Figure 1. Graph of the Dirichlet function.
In [14] another approach has been developed. The introduction of [14] reads: “after
all, analysis has few bases: the concepts of function, infinity, limit, continuity, and
differentiability. To visit these locations in a comprehensive way, one should allow
these concepts the faculty of expressing all of their potential and all their fantastic
and incredible situations. To achieve this high level it is necessary to stick closely to
the definition of function as a correspondence of free sets.” The modern definition
of function clearly allows a deeper comprehension of other concepts, such as those
of limit, continuity, and differentiability. It is important that students understand the
far-reaching features of these notions, because only in this way will they be able to
avoid mistakes due to a limited concept of function. Students should be aware from
the very beginning that formulas provide a very small class of functions: analyticity is
a property enjoyed by very special functions. Being discontinuous is not synonymous
with being pathological because most functions are discontinuous. In the same way,
being continuous is not synonymous with being smooth since continuous functions
can be very rough, as in the case of the celebrated Weierstrass functions.
Following this plan, in [14] the authors propose a collection of problems the solu-
tions of which require the modern definition of function: in some cases, using highly
discontinuous functions is not essential but it simplifies the argument and makes the
situation clear.
We note that there is a vast literature devoted to so-called pathological functions. We
mention the classic book by Gelbaum and Olmsted [6] and the recent extensive mono-
graph by Kharazishvili [7]. We also mention Thim [12], which is basically a treatise
concerning continuous nowhere differentiable functions. However, at a pedagogical
level, not much material seems to be available.
In this paper, we adopt the approach of [14] and we discuss some peculiar examples,
some of which are taken from [14]. The focus is on special discontinuous functions,
the definitions of which are algebraically and technically simple in most cases. We
believe that this approach could be used not only with good students but also with weak

students who may take advantage of examples which illustrate profound concepts but
remain at a low level of complexity.
2. EXAMPLES OF DIRICHLET-TYPE FUNCTIONS. In this section we present

several discontinuous functions modeled on the Dirichlet function (1). We think that
students may find such examples easy and eventually amusing, once they are ac-
quainted with function (1). On the other hand, asking a student of a first-year calculus
course to provide such examples could be very challenging.
1. A function continuous at only one point. If we ask a beginner in mathematical

analysis whether a function from R to itself can be continuous at only one point, most
of the time we will get the answer “no” as a consequence of an over-simplified view
of the concept of continuity. Usually, students have an idea of continuity as a global
property of a function: according to this point of view, continuous functions are those
functions whose graphs can be drawn without lifting the chalk from the blackboard.
This vision is not completely wrong since a function is continuous in the whole of an
interval if and only if its graph is a pathwise-connected set in the plane. However, the
following example clarifies the local nature of continuity.
–3 –2 –1 1 2 3
–1
–2
–3
Figure 2. Graph of function F1 .
Let F1 be the function from R to itself defined by

(
x, if x ∈ Q,
F1 (x) =
−x, if x ∈ R \ Q.
The function F1 is continuous only at x = 0.

As a variant of the function F1 , one can consider the function G 1 from R to itself
defined by
(
sin x, if x ∈ Q,
G 1 (x) =
− sin x, if x ∈ R \ Q.

2
–2 – 2
–1
–2
Figure 3. Graph of function G 1 .
The function G 1 is continuous only at a countable set of points, precisely at the points
x = kπ with k ∈ Z.
2. A function differentiable at only one point. Loosely speaking, the derivative of

a function at a point is the slope of the tangent line to the graph of the function at
that point. However, the notion of differentiability has a local nature which goes much
beyond the geometric idea of a tangent line.
Let F2 be the function from R to itself defined by
(
x 2, if x ∈ Q,
F2 (x) = 2
−x , if x ∈ R \ Q.
The function F2 is differentiable only at x = 0 and is discontinuous at any point

x 6 = 0. Proving that F2 is differentiable at x = 0 is an easy exercise. However, it may
be instructive to deduce such a proof from the following elementary lemma.
Lemma 2. Let f, g, and h be functions from R to itself such that f (x) ≤ h(x) ≤ g(x)
for all x ∈ R. Let x0 ∈ R be fixed. If f and g are differentiable at x0 , f ′ (x0 ) = g ′ (x0 ),
and f (x0 ) = g(x0 ), then h is differentiable at x0 and h ′ (x0 ) = f ′ (x0 ) = g ′ (x0 ).
3 3
2 2
1 1
–2 –1 1 2 –2 –1 1 2
–1 –1
–2 –2
–3 –3
Figure 4. Graph of function F2 . Figure 5. A randomly generated function differ-
entiable at x = 0.

Students may find it entertaining to apply the previous lemma to randomly gener-
ated functions h satisfying the condition −x 2 ≤ h(x) ≤ x 2 for all x in a neighborhood
of zero and find out that such functions are differentiable at zero (see Figure 5).
3. A function with positive derivative at one point, which is not monotone in any
neighborhood of that point. Let F3 be the function from R to itself defined by
(
x 2, if x ∈ Q,
F3 (x) =
2x − 1, if x ∈ R \ Q.
By Lemma 2, the function F3 is differentiable at the point x = 1 and F3′ (1) = 2, but
F3 is not monotone in any neighborhood of 1. Those students who are used to identify-
ing increasing functions with functions with positive derivative may find this example
rather bizarre. However, this example points out a well-known subtle distinction con-
cerning increasing functions. In fact, if a function f from R to itself is differentiable
at a point x0 and f ′ (x0 ) > 0 then f is increasing at the point x0 in the sense that
for all x in a convenient neighborhood of x0 we have: f (x) > f (x0 ) if x > x0 and
f (x) < f (x0 ) if x < x0 . Monotonicity is a stronger notion and, in the case of differ-
entiable functions, it occurs when the derivative does not change sign in the whole of
an interval.
Variants of the function F3 can be easily defined by considering convex combina-

tions of x 2 and 2x − 1 of the form
θ (x)x 2 + (1 − θ (x))(2x − 1),
where θ is a suitable function taking values in [0, 1]. A “continuous version” of F3 is

given for example by the function G 3 from R to itself defined by1

2 9 2 2 9
G 3 (x) = sin · x + 1 − sin · (2x − 1),
x −1 x −1
1 The choice of the number 9 in the definition of G 3 is not essential, and is only aimed at emphasizing the
oscillations of the graph in Figure 7.

for all x ∈ R \ {1}, and G 3 (1) = 1. Again, G ′3 (1) = 2 > 0, but G 3 is not monotone in
any neighborhood of the point x = 1.
Clearly the function G 3 is much smoother than F3 ; however, defining F3 and prov-
ing its nonmonotonicity is simpler.
2.5
2.0
1.5
1.0
0.5
–0.5 0.5 1.0 1.5
–0.5
3. FURTHER EXAMPLES. In this section we present examples of functions with a

specific type of discontinuity at any point of the domain, such as removable discontinu-
ities or jump discontinuities. Although complicated, the definitions of these functions
are constructive and make use of elementary notions. We note that imposing a specific
type of discontinuity at any point of the domain of a function leads to limitations on
the structure of the domain itself. In fact, the following result in Klippert [10] holds.
Theorem 3. Let f be a function from R to itself. Let A be the set of those points x0 ∈ R
such that f is discontinuous at x0 and at least one of the two limits limx→x + f (x) and
0
limx→x − f (x) exists and is finite. Then A is at most countable.
0
4. A function with removable discontinuities at every point of its domain. The

following function F4 is defined on the dyadic rational numbers in (0, 1), i.e., those
numbers in (0, 1) with a finite binary expansion. A simple formal definition is given
below in (6). However, we prefer to begin by describing the elementary geometric
construction given in [14], which could also be presented to beginners.
We construct the function F4 by means of the following iterative procedure (see
Figure 9).
Step 1. Let V = 0 and L = 1, and consider the square in the Euclidean plane with
edges parallel to the coordinate axes whose lower left vertex is (V, V ) and

1.0
0.8
0.6
0.4
0.2
0.0
0.0 0.2 0.4 0.6 0.8 1.0
Figure 8. Graph of function F4 (arrested at the fourth iteration).
whose edges have length L. Let

n
X L
BV,L = V+ : n∈N .
j=1
2j
Then we set F4 (x) = V + L for all x ∈ BV,L .

Step 2. Inside the square defined in Step 1, consider the sequence of squares with
edges of length L n = 2Ln and parallel to the coordinate axes, and lower left
vertices with coordinates (Vn , Vn ), where Vn = V + n−1 L
P
j=1 2 j for all n ∈ N
with n ≥ 2, and (V1 , V1 ) = (V, V ).
Repeat Step 1 and Step 2, using each square obtained in Step 2 in place of
the original square, and iterate the process.
The domain of the function F4 defined by this procedure is B = ∪V,L BV,L , which is
clearly the set of the dyadic rational numbers in (0, 1).
Looking closely at the construction described above and at the few first iterations
(see Figure 9), it is evident that
lim F4 (x) = v < F4 (v) (4)

x→v
for all v ∈ B. Thus F4 has a removable discontinuity at each point v ∈ B. Indeed, in

order to remove the discontinuity at a point v, it is enough to redefine the function by
setting F̃4 (v) = v.
It is possible to give a non-iterative definition of the function F4 . Indeed the set B
can be represented as

1.0 1.0
0.8 0.8
0.6 0.6
0.4 0.4
0.2 0.2
0.0 0.0
0.0 0.2 0.4 0.6 0.8 1.0 0.0 0.2 0.4 0.6 0.8 1.0
1.0 1.0
0.8 0.8
0.6 0.6
0.4 0.4
0.2 0.2
0.0 0.0
0.0 0.2 0.4 0.6 0.8 1.0 0.0 0.2 0.4 0.6 0.8 1.0
Figure 9. The first four iterations.
np o
n
B= : n, p ∈ N, p is odd, p < 2 (5)
2n
and F4 can be defined on B directly by the equality
p p+1
F4 n = (6)
2 2n
for all n, p ∈ N, where p is odd and p < 2n . This alternative definition allows an
easy proof of (4) since the summand 1/2n in the right-hand side of (6) vanishes in the
limiting procedure (see also (8) and (9)).
A well-known function which enjoys a similar property can be found in several
textbooks. It is defined as follows. Let G 4 be the function from (0, 1) ∩ Q to itself
defined by

p 1
G4 = (7)
q q

1
2
1
3
1
4
1
5
1 1 1 2 3
4 3 2 3 4
1
for all p, q ∈ N, 0 < p/q < 1, p, q coprime. It is not difficult to prove that
lim G 4 (x) = 0 6= G 4 (x0 ) (8)

x→x0
for all x0 ∈ (0, 1) ∩ Q. The function G 4 can be extended to the whole of (0, 1) by
setting G 4 (x) = 0 for all x ∈ (0, 1) \ Q and turns out to be continuous at any point
x ∈ (0, 1) \ Q. According to Appell [1, p. 10] some students would call function G 4
“Dirichlet light” because its set of discontinuities is “lighter” than the set of disconti-
nuities of the Dirichlet function.
We note that by (6) and (7) it immediately follows that
F4 (x) = x + G 4 (x) (9)
for all x ∈ B. Equality (9) could be used to define F4 via G 4 and to give a further
proof of (4). In this way, the function F4 , whose original definition is based on a geo-
metric construction, can eventually be represented via a function whose definition has
a completely different origin.
5. A function with jump discontinuities at all points of its domain. Monotone func-
tions with many jump discontinuities are well known in real analysis and probability
theory: these functions can be obtained as cumulative distribution functions associated
with discrete measures (see, e.g., Folland [5, p. 102]). Here we present two functions
defined on the set B of the dyadic rational numbers in (0, 1) with jump discontinuities
at every point of B. These functions are defined directly by using binary notation and
do not require any advanced tool.
We recall that the set B in (5) can be represented also in the form
∞ Xn
[ x j
B= j j
x ∈ {0, 1} for all j ∈ {1, . . . , n}, xn = 1 . (10)
n=1 j=1
2
Thus, given x ∈ B there exists n(x) ∈ N such that x = n(x) j

P
j=1 x j /2 for suitable values
of x j ∈ {0, 1}, where xn(x) = 1; using binary notation x can be written as
0. x1 x2 · · · xn(x) .
Let F5 be the function from B to itself defined by
F5 (0. x1 x2 · · · xn(x) ) = 0. y1 y2 · · · ym(x) ,

where the digits y1 , . . . , ym(x) are obtained from the digits x1 , . . . , xn(x) of x by
simply replacing each 0 with 01. For example, F5 (0.01) = 0.011, F5 (0.101001) =
0.101101011, F5 (0.000011) = 0.0101010111, and F5 (0.1) = 0.1. It is easy to prove
that F5 is a monotone function strictly increasing on B. Thus, F5 admits left and right
limits at any point of B. Moreover, it is not difficult to realize that in fact
F5 (v) = lim F5 (x) < lim F5 (x)

x→v − x→v +
for all v ∈ B. Hence F5 is left-continuous with a jump discontinuity at every point

v ∈ B. In order to get rid of the monotonicity one can consider the function G 5 defined
by
G 5 = F4 ◦ F5 .
The function G 5 is particularly interesting since it admits a geometric description in the

spirit of the iterative construction of F4 : in this case one has to consider rectangles with
no common vertices rather than squares (see Figures 11 and 12). The function G 5 is
discussed in detail in Drago [4], where it is also proved by using binary representations
that for all v ∈ B the one-sided limits of G 5 at v exist and satisfy
lim G 5 (x) < lim G 5 (x) < G 5 (v).

x→v − x→v +
1.0
0.8
0.6
0.4
0.2
0.0
0.0 0.2 0.4 0.6 0.8 1.0
Figure 11. Graph of function G 5 (arrested at the third iteration).
6. A function approaching infinity at any point of its domain. The following ex-
ample is a simple variant of the function G 4 . Let F6 be the function from (0, 1) ∩ Q to
itself defined by

p
F6 =q
q

1.0
0.8
0.6
0.4
0.2
0.0
0.0 0.2 0.4 0.6 0.8 1.0
Figure 12. Construction of function G 5 .
50
40
30
20
10
0
1 1 1 2 3 1
4 3 2 3 4

for all p, q ∈ N, 0 < p/q < 1, p, q coprime. It is not difficult to prove that
lim F6 (x) = +∞
x→x0
for all x0 ∈ [0, 1].
ACKNOWLEDGMENTS. The authors are thankful to Prof. Maurizio Emaldi for pointing out reference [9]
and useful discussions, and also to Dr. Matteo Dalla Riva for pointing out reference [10]. The first two authors
are deeply indebted to the third author, P. Toni, who, in different periods, was their teacher of mathematics
and physics in secondary school, where he has developed and experienced the main pedagogical ideas that
are the basis of the present paper for about three decades. Those stimulating lessons led the second author,
P. D. Lamberti, to define the function F4 . The functions F5 and G 5 were defined by the first author, G. Drago,
in his dissertation [4], written under the supervision of the second author.
The authors are very thankful to a referee for suggesting formula (6), which simplifies a binary representa-
tion given in the first version of this paper, and for pointing out formula (9).
REFERENCES
1. J. Appell, Analysis in Beispielen und Gegenbeispielen, Springer, Berlin, 2009.

2. N. Bourbaki, Eléments de Mathématique, Théorie des Ensémbles, Hermann, Paris, 1970.
3. R. B. Deal, Some remarks on “A dilemma in definition,” Amer. Math. Monthly 74 (1967) 1258–1259.
http://dx.doi.org/10.2307/2315696
4. G. Drago, Funzioni e Funzioni Strane: Alcuni Aspetti Storici e Didattici, B.A. thesis, University of
Padova, Padova, 2009.
5. G. B. Folland, Real Analysis: Modern Techniques and Their Applications, John Wiley, New York, 1984.
6. B. R. Gelbaum and J. M. H. Olmsted, Counterexamples in Analysis, Holden-Day, San Francisco, 1964.
7. A. B. Kharazishvili, Strange Functions in Real Analysis, 2nd ed., Pure and Applied Mathematics, vol.
272, Chapman & Hall/CRC, Boca Raton, FL, 2006.
8. I. Kleiner, Evolution of the function concept: A brief survey, College Math. J. 20 (1989) 282–300. http:
//dx.doi.org/10.2307/2686848
9. , Functions: Historical and pedagogical aspects, Science & Education 2 (1993) 183–209. http:
//dx.doi.org/10.1007/BF00592206
10. J. Klippert, Advanced advanced calculus: Counting the discontinuities of a real-valued function with
interval domain, Math. Mag. 62 (1989) 43–48. http://dx.doi.org/10.2307/2689945
11. C. P. Nicholas, A dilemma in definition, Amer. Math. Monthly 73 (1966) 762–768. http://dx.doi.
org/10.2307/2313993
12. J. Thim, Continuous Nowhere Differentiable Functions, M.A. thesis, Luleå University of Technology,
Luleå, Sweden, 2003.
13. H. A. Thurston, A reply to “A dilemma in definition,” Amer. Math. Monthly 74 (1967) 1259–1261. http:
//dx.doi.org/10.2307/2315697
14. P. Toni and P. D. Lamberti, Esplorando l’Analisi Matematica, Società Editrice Internazionale, Torino,
1996.
15. A. P. Youschkevitch, The concept of function up to the middle of the 19th century, Arch. Hist. Exact Sci.
16 (1976/77) 37–85.
GIACOMO DRAGO is 22 years old. He got his bachelor’s degree with honors from the University of Padova
in 2009. He distinguished himself in the national finals of three consecutive editions of the Kangourou Mathe-
matics Competitions in Italy (2004–2006). He is currently working on his master’s degree in mathematics, and
he dreams of becoming a teacher of mathematics.
Dipartimento di Matematica Pura ed Applicata, Università degli Studi di Padova, Via Trieste 63, 35121
Padova, Italy
giacomo.drago@gmail.com
PIER DOMENICO LAMBERTI received his Ph.D. in mathematics from the University of Padova in 2003.
He currently has a position as researcher in mathematical analysis at the same university. His research in-

terests are mainly in spectral stability problems for elliptic partial differential equations. He likes to think of
mathematics as a form of art.
Dipartimento di Matematica Pura ed Applicata, Università degli Studi di Padova, Via Trieste 63, 35121
Padova, Italy
lamberti@math.unipd.it
PAOLO TONI received his master’s degree in mathematics from the University of Firenze in 1972. He has
taught mathematics and physics in high school for almost 40 years. He was a pioneer in the organization of
mathematics competitions in Italy. His interests are mainly in mathematical creativity and teacher education.
He retired in 2010. Now, he enjoys staying with his nephews and playing the organ. He believes in students’
creativity.
Liceo Scientifico Statale E. Fermi, Via Vittorio Emanuele II 50, 35100 Padova, Italy
paolotoni@inwind.it
n+ 12
1
A Ramanujan-Type Formula for 1+ e
n
n+ 12
Motivated by Sanjay Khattri’s proofs [1] that e < 1 + n1
!
, I offer the follow-
ing formula, which I think Ramanujan would have liked, for the ratio of the two
sides of Khattri’s inequality:
1
1 n+ 2
Z 1
u2

1+ e = exp 2 2
du
n 0 (2n + 1) − u
Z 1
u2

= exp 2
du (where m = n(n + 1)/2)
0 8m + 1 − u
Z 1
αu 2

= exp 2
du (where α = 1/(8m))
0 1 + α(1 − u )
Z 1
2 2 2 2 3 2 2 2 4 2 2 3
= exp αu − α u (1 − u ) + α u (1 − u ) − α u (1 − u ) + · · · du
0

1 1 2 2 1 4 2 3 1 6 4 2 4 1 8 6 4 2 5
= exp α− · α + · · α − · · · α + · · · · α − ···
3 5 3 7 5 3 9 7 5 3 11 9 7 5 3

α 2α 4α 6α 8α !
= exp 1− 1− 1− 1− 1 − ···
3 5 7 9 11

1 1 2 3 4 !
= exp 1− 1− 1− 1− 1 − ··· .
24m 5 · 4m 7 · 4m 9 · 4m 11 · 4m
REFERENCES
n+0.5
1. S. K. Khattri, Three proofs of the inequality e < 1 + n1 , Amer. Math. Monthly 117 (2010)
273–277. http://dx.doi.org/10.4169/000298910X480126
—Submitted by Michael D. Hirschhorn,

University of New South Wales, Sydney, Australia

An Introduction to Gauss Factorials
John B. Cosgrave and Karl Dilcher
Abstract. Starting with Wilson’s theorem and its generalization by Gauss, we define a Gauss
factorial Nn ! to be the product of all positive integers up to N that are relatively prime to n. We
present results on the Gauss factorials ( n−1
M n
) !, and more generally on partial products obtained
when the product (n − 1)n ! is divided into M equal parts, for integers M ≥ 2. Finally, exten-
sions of the Gauss binomial coefficient theorem are presented in terms of Gauss factorials.
1. INTRODUCTION. One of the most remarkable results in elementary number the-

ory is Wilson’s theorem and its converse by Lagrange, stating that p is a prime if and
only if
( p − 1)! ≡ −1 (mod p). (1.1)
A proof of this result can be found in most introductory books on number theory, and
it depends on the fact that if p is prime then any integer a with 1 < a < p − 1 has its
inverse a −1 6 ≡ a (mod p).
For any odd prime p, if we write out the factorial ( p − 1)! and exploit symmetry
modulo p, we obtain
p−1 p+1
! p−1 p−1 !
p−1

1 · 2··· 2 2
· · · ( p − 1) ≡ 2
!(−1) 2
2
! (mod p), (1.2)
and therefore, by (1.1),

! p−1 2 p+1
2
! ≡ (−1) 2 (mod p). (1.3)
This was apparently first observed by Lagrange (see [13, p. 275]), and this congru-
ence can be used, along with a result of Mordell [24] that involves the class numbers
of imaginary
! p−1 quadratic fields, to completely characterize the multiplicative order of
2
! modulo p. We will leave this aside, and instead consider now the two halves
of the product on the left-hand side of (1.2). We denote these two partial products by
5(2) (2)
1 and 52 , respectively, where the upper index indicates the fact that we divide
the entire product into two equal parts. Using Wilson’s theorem (1.1) and symmetry
modulo p, we obtain 5(2) (2)
1 52 ≡ −1 (mod p) and
p−1
52(2) ≡ (−1) 2 5
(2)
1 (mod p). (1.4)
This is, of course, equivalent to (1.3), but writing it in this way gives rise to the follow-
ing question:
What can we say about the three partial products 5(3) (3) (3)
1 , 52 , 53 obtained by di-
viding the entire product ( p − 1)!, that is, the left-hand side of (1.1), into three equal
parts? For this we require p to be of the form p ≡ 1 (mod 3) or, in fact (since p is
prime), p ≡ 1 (mod 6); we then have
5(3)
1 = 1 · 2···
p−1
3
, 5(3)
2 =
p+2
3
· · · 2 p−2
3
, 5(3)
3 =
2 p+1
3
· · · ( p − 1).

In analogy to (1.4) we see an obvious symmetry relation between 5(3) (3)
1 and 53 ,
namely
5(3) (3)
3 ≡ 51 (mod p), (1.5)
but without a power of −1 since p−13

is always even. However, there is no obvious
relation between 51 and the “middle third” 5(3)
(3)
2 .
Going one step further, for p ≡ 1 (mod 4) we now divide the entire product ( p −
1)! into four partial products
! ! ! p−1
5(4)
j = ( j − 1) p−1
4
+ 1 ( j − 1) p−1
4
+ 2 · · · j 4 , ( j = 1, 2, 3, 4). (1.6)
This time we have two obvious symmetry relations, namely

p−1 p−1
5(4)
4 ≡ (−1)
4 5
(4)
1 (mod p), 5(4)
3 ≡ (−1)
4 5
(4)
2 (mod p), (1.7)
while there is no obvious relation between 5(4) 1 and 52 .

(4)
Table 1 illustrates the congruences (1.5) and (1.7). We also see that, indeed, there
are no obvious relationships between 5(M) 1 and 52(M) for M = 3 and 4 and p < 100,
with the exceptions of p = 7 and p = 61, where 5(3) (3)
1 ≡ −52 (mod p). These could,
of course, be coincidences, but it turns out that this congruence holds also for p = 331,
p = 547, p = 1951, and for further relatively rare primes, as explained later. Here it
is interesting to mention p = 3571, the first case with 5(3) (3)
1 ≡ −52 ≡ 1 (mod p). In
contrast, there are no primes p for which 5(3) (3)
1 ≡ 52 (mod p), p ≡ 1 (mod 6), or
5(4) (4)
1 ≡ ±52 (mod p), p ≡ 1 (mod 4). All this is readily explained by appealing to
two deep theorems of Jacobi and of Gauss; we will return to this later.
Table 1. The partial products modulo p, for ( p − 1)! split into M = 3 and M = 4 equal parts.
p 5(3)
1 5(3)
2 5(3)
3 p 5(4)
1 5(4)
2 53(4) 54(4)
7 2 −2 2 5 1 2 −2 −1
13 −2 3 −2 13 6 3 −3 −6
19 −2 −5 −2 17 7 −3 −3 7
31 2 −8 2 29 −6 −2 2 6
37 7 3 7 37 −16 5 −5 16
43 −3 19 −3 41 13 7 7 13
61 −14 14 −14 53 26 7 −7 −26
67 −20 −33 −20 61 19 7 −7 −19
73 33 −12 33 73 18 −35 −35 18
79 −37 3 −37 89 22 42 42 22
97 21 −11 21 97 20 −28 −28 20
This naturally leads to the question of dividing the product ( p − 1)! into 5, 6, or
in general M ≥ 2 partial products of equal length, for primes p ≡ 1 (mod M). In
analogy to (1.6) we define, for such M and p, the products
p−1
Y
M
!
5(M)
j = ( j − 1) p−1
M
+i , ( j = 1, 2, . . . , M). (1.8)
i=1
November 2011] GAUSS FACTORIALS 813

Once again it is clear that
5(M) (M)
M− j ≡ ±5 j (mod p), j = 1, 2, . . . , ⌊ M−1
2
⌋,
(M)
with the “central product” 5(M+1)/2 playing a somewhat special role when M is odd.
Extensive computations suggest that while (for fixed M) there are instances where two
of the partial products, with 1 ≤ j ≤ ⌊ M+1
2
⌋, are congruent, there are no cases where
all are simultaneously congruent.
2. COMPOSITE MODULI. This might well have been the end of the story were it
not for the possibility of considering composite moduli. Since our point of departure
has been Wilson’s theorem (1.1), let us first recall why Lagrange’s converse is true. If
n is composite, we can write it as n = n 1 n 2 , with 1 < n 1 < n. But then n 1 | (n − 1)!,
and therefore (n − 1)! 6 ≡ ±1 (mod n). However, if we suitably modify the factorial
on the left-hand side of (1.1), we obtain a composite analogue of Wilson’s theorem. It
was Gauss who first proved the following theorem.
Theorem 1 (Gauss). For any integer n ≥ 2 we have

(
Y −1 (mod n) for n = 2, 4, p α , or 2 p α ,
j≡ (2.1)
1≤ j≤n−1
1 (mod n) otherwise,
gcd( j,n)=1
where p is an odd prime and α is a positive integer.
The number of integers j in the product in (2.1), that is, those positive integers up
to n that are relatively prime to n,Qis given by Euler’s totient function φ(n), which has
the explicit evaluation φ(n) = n p|n (1 − 1p ), with the product taken over all prime
divisors p of n. We recall that the integers n in (2.1) for which the product is −1
(mod n) are exactly those that have a primitive root; this fact is important for the proof
of the result.
In spite of the fact that Theorem 1 was stated in the famous Disquisitiones Arith-
meticae [17, §78] and in the equally influential books [14, §38] and [20, p. 102], sur-
prisingly little can be found on this topic in the literature. The few published references
to this result include [21] and [27], where Theorem 1 was further extended, and [22]
and [1], where (2.1) was used to extend the classical Wilson quotient to composite
moduli. The theorem was rediscovered at least once; see [26].
In order to state this theorem and numerous other results more concisely, we intro-
duce the following notation: for positive integers N and n let Nn ! denote the product
of all integers up to N that are relatively prime to n, i.e.,
Y
Nn ! = j. (2.2)
1≤ j≤N
gcd( j,n)=1
This notation is a slight variation of the one used in [18], a useful reference on factorial
and binomial congruences. To be able to refer more easily to Nn !, we shall call it a
Gauss factorial, a terminology suggested by Theorem 1, which we call from here on
the Gauss-Wilson theorem.
We now turn to the composite analogue of Lagrange’s observation! n−1 in (1.2) and (1.3)
and begin with a general discussion of the Gauss factorial 2 n ! for odd integers

n ≥ 3. Since (n − 1)n ! is a product of φ(n) residues, and φ(n) is even for odd n ≥ 3,
then by the same symmetry argument as in (1.2) we obtain
! n−1 2 1
2 n
! ≡ (−1) 2 φ(n)+ε (mod n), (2.3)
where, by (2.1), ε = 1 when n = p α , and ε = 0 otherwise. Now φ( p α ) =

( p − 1) p α−1 , and therefore
p−1 p+1
1
2
φ( p α ) +1≡ 2
+1= 2
(mod 2).
On the other hand, φ(n) is divisible by 4 if n has at least two distinct odd prime factors.
Hence by (2.3) we get
(
! n−1 2 −1 (mod n) if n = p α , p ≡ 1 (mod 4),
2 n
! ≡ (2.4)
1 (mod n) otherwise.
This is analogous to (1.3) for prime powers. In connection with this congruence we
remark that the multiplicative orders of ( n−1
2 n
) ! modulo n were completely determined
in [7].
As we did following (1.3), we will now turn to dividing the product (n − 1)n ! into
partial products. In complete analogy to (1.8) we define our partial products 5(M)
j , for
integers M ≥ 2 and n ≡ 1 (mod M), as
Y
5(M)
j := i, ( j = 1, 2, . . . , M), (2.5)
(M)
i∈I j
where, for j = 1, 2, . . . , M,

I j(M) := i | ( j − 1) n−1
M
+ 1 ≤ i ≤ j n−1
M
, gcd(i, n) = 1 . (2.6)
The dependence on n, always considered a fixed modulus, is implied in this notation.

When n = p is prime, the definition (2.5) reduces to (1.8); the use of identical notation
is therefore justified.
Table 1 could now be extended to include composite moduli. To save space, only the
last ten cases with n < 100 for each of M = 3 and M = 4 are shown in Table 2. This
immediately shows that, in contrast to the prime modulus case, the partial products can
indeed all be congruent to each other modulo n. In particular, we see that this happens
for n = 91 when M = 3, and for n = 65 and 85 when M = 4.
Table 2. Partial products modulo n, for 10 values of n < 100, M = 3, 4.
n 5(3)
1 5(3)
2 5(3)
3 n 5(4)
1 5(4)
2 5(4)
3 5(4)
4
70 29 1 29 61 19 7 −7 −19
73 33 −12 33 65 8 8 8 8
76 −29 −15 −29 69 31 −26 −26 31
79 −37 3 −37 73 18 −35 −35 18
82 −33 −25 −33 77 16 31 31 16
85 −28 9 −28 81 2 40 −40 2
88 5 −7 5 85 13 13 13 13
91 29 29 29 89 22 42 42 22
94 −23 43 −23 93 34 −10 −10 34
97 21 −11 21 97 20 −28 −28 20

3. A QUESTION CONSIDERED BY D. H. LEHMER. Before explaining this last
observation in Section 4, let us pause briefly to consider the sets I j(M) defined in (2.6),
and in particular their cardinalities
φ M, j (n) := #I j(M) . (3.1)
First, when n = p is a prime, then clearly for a given M all I j(M) have the same number
of elements, namely φ M, j ( p) = p−1M
. Next, when M = 1, then φ1,1 (n) = φ(n). When
M = 2, then by symmetry of the sets I1(2) and I2(2) we have φ2,1 (n) = φ2,2 (n) = 12 φ(n).
However, already in the case M = 3 the situation is not as straightforward, as the
example n = 4 shows: in this case we have φ3,1 (n) = φ3,3 (n) = 1, but φ3,2 (n) = 0.
According to D. H. Lehmer [23] it was J. J. Sylvester who coined the term totatives
for those positive integers up to a given n that are relatively prime to n. We are therefore
dealing with the distribution of totatives in subintervals of the interval [1, n]. Lehmer
[23] was the first to study this distribution, and to give a sufficient condition for the
equal distribution of totatives. This area of study has attracted the attention of later
mathematicians; for instance, a conjecture of Erdős [16] was proven by Hall and Shiu
[19].
Table 3. The first ten moduli n for which all φ M, j (n) are equal, for each of M = 3, 4, 5.
n factored φ3, j n factored φ4, j n factored φ5, j
28 22 · 7 4 25 52 5 66 2 · 3 · 11 4
49 72 14 45 2
3 ·5 6 121 112 22
2 4
52 2 · 13 8 65 5 · 13 12 176 2 · 11 16
70 2·5·7 8 85 5 · 17 16 186 2 · 3 · 31 12
76 22 · 19 12 105 3·5·7 12 231 3 · 7 · 11 24
91 7 · 13 24 117 32 · 13 18 246 2 · 3 · 41 16
112 24 · 7 16 125 53 25 286 2 · 11 · 13 24
2
124 2 · 31 20 145 5 · 29 28 341 11 · 31 60
130 2 · 5 · 13 16 153 32 · 17 24 366 2 · 3 · 61 24
133 7 · 19 36 165 3 · 5 · 11 20 396 22 · 32 · 11 24
Table 3 seems to indicate that whenever n ≡ 1 (mod M) has a prime factor p sat-
isfying p ≡ 1 (mod M), then the corresponding totatives are equally distributed, that
is, all φ M, j (n) are equal. This is in fact true, as was shown by D. H. Lehmer [23,
Theorem 4]:
Lemma 1 (Lehmer). Let M ≥ 2 and n ≡ 1 (mod M). If n has at least one prime
factor p with p ≡ 1 (mod M), then the totatives in the interval [1, n] are equally
distributed, that is, we have
1
φ M, j (n) = φ(n), ( j = 1, 2, . . . , M). (3.2)
M
Lehmer actually showed something slightly different, namely that under the given
conditions the intervals
n n
( j − 1) <i < j , j = 1, 2, . . . , M, (3.3)
M M

have equal numbers of integers i relatively prime to n, and the endpoints cannot them-
selves be totatives. However, we can easily see that this implies Lemma 1. For further
details, see [7, Lemma 2]. This lemma is an important ingredient in a proof in the next
section.
While Lehmer’s theorem gives a sufficient condition, the following example shows
that it is not necessary: Let M = 8 and n = 105 = 3 · 5 · 7. Although none of the
prime factors of n are of the form p ≡ 1 (mod 8), a computation shows that each of
the eight sets I j(8) contains 18 φ(105) = 6 elements.
4. THE CASE WHERE THE PARTIAL PRODUCTS ARE ALL CONGRUENT.

In order to further explore the case where all partial products 5(M)
j , for a fixed M, are
congruent to each other modulo n, we computed many pairs of M and n for which
this is the case. Table 4 shows the first ten such moduli n for each of M = 3, 4, and 5,
along with the factorizations of n and the common values
! n−1
M n
! ≡ 5(M)
j (mod n), j = 1, 2, . . . , M. (4.1)
By definition the left-hand side of (4.1) is obviously identical with the right-hand side
for j = 1.
(M)
Table 4. The first ten moduli n for which all 5 j are congruent, for each of M = 3, 4, 5.
Q(5)
n factored 5(3)
j n factored 5(4)
j n factored j
91 7 · 13 29 65 5 · 13 8 341 11 · 31 −85
133 7 · 19 58 85 5 · 17 13 451 11 · 41 −105
217 7 · 31 67 145 5 · 29 1 671 11 · 61 −304
244 22 · 61 1 185 5 · 37 −68 781 11 · 71 −117
247 13 · 19 −88 205 5 · 41 1 1111 11 · 101 36
259 7 · 37 100 221 13 · 17 −1 1271 31 · 41 264
301 7 · 43 36 265 5 · 53 23 1441 11 · 131 89
364 22 · 7 · 13 113 305 5 · 61 −121 1661 11 · 151 −545
403 13 · 31 118 325 52 · 13 −57 1891 31 · 61 497
427 7 · 61 135 365 5 · 73 27 1991 11 · 181 125
We see from this table that all the moduli, except n = 244 = 22 · 61, have at least
two distinct prime factors that are congruent to 1 modulo M. In fact, we have:
Theorem 2. Let M ≥ 2 be an integer, and suppose that the positive integer n has at
least two distinct prime factors congruent to 1 (mod M). Then all the partial products
5(M)
j are congruent modulo n, that is, the congruences (4.1) hold.
Our starting point for the proof of Theorem 2 is the observation that each partial
product 5(M)
j can be written as a quotient of two Gauss factorials that are similar in
nature. In particular, we see immediately from the definitions (2.5) and (2.2) that
! n−1
(M) j M n!
5j = ! , j = 1, 2, . . . , M, (4.2)
( j − 1) n−1
M n
!
with the convention that 0n ! = 1. We therefore need to study the Gauss factorials on the
right-hand side of (4.2). Our main tool is the following explicit congruence, obtained
as a generalization of Proposition 2 in [7].

Lemma 2. Let M ≥ 2 and n ≡ 1 (mod M), n = p α q β w for distinct prime p, q ≡ 1
(mod M), α, β ≥ 1, and gcd( pq, w) = 1. Then for i = 1, 2, . . . , M we have
p−1
! n−1 εi M p α−1
i M n ! ≡ i A (mod q β w), A= φ(q β w), (4.3)
p M
where ε = −1 if w = 1, and ε = 1 if w > 1, and φ(m) denotes Euler’s totient func-

tion.
Now, combining the congruence (4.3) with (4.2), we get

p−1
ε M p α−1
5(M)
j ≡ (mod q β w), A= φ(q β w). (4.4)
pA M
Since p α and q β are interchangeable, we also have

p−1
ε M q β−1
5(M)
j ≡ B (mod p α w), B= φ( p α w). (4.5)
q M
By the Chinese remainder theorem, applied to (4.4) and (4.5), the partial product 5(M)j
is uniquely determined modulo p α q β w = n, and most importantly, it is independent
of j. This completes the proof of Theorem 2.
To prove Lemma ! 2,
the main idea is to break the whole range of the product in the
Gauss factorial i n−1
M n
! into a number of products of approximately equal length and
a shorter “tail.” We then evaluate the products of the first type using the Gauss-Wilson
theorem with modulus e n := q β w. To do so we divide i n−1
M
by e
n with remainder:
n−1 e
n−1 pα − 1
i = ise
n+i , where s := . (4.6)
M M M
By hypothesis we know that s and (en − 1)/M are both integers. Based on (4.6) we
now decompose our Gauss factorial into is products of similar lengths and one shorter
product; that is, we write
 
is
Y
! n−1
i M n! =  P j  Q, (4.7)
j=1
where
e
n −1 ie
n −1
Y ! Y M
!
P j := ( j − 1)e
n+k , Q := ise
n+k . (4.8)
k=1 k=1
gcd(( j−1)e
n +k,n)=1 gcd(ise
n +k,n)=1
For a given j, if the set of residues {( j − 1)e n+k | 1 ≤ k ≤e n − 1}, subject to gcd(( j −
n + k, n) = 1, formed a reduced residue system modulo e
1)e n , then the product P j
would, by the Gauss-Wilson theorem, be congruent to −1 (mod e n ) if w = 1, and to
1 (mod e n ) if w > 1. However,
! n−1 this is not always the case because the residues that
appear in the product i M n ! have none divisible by p; these residues have been re-
!
moved from the normal factorial i n−1 M
! in forming the corresponding Gauss factorial.

To deal with the variable nature of these P j , we multiply all relevant multiples of
p back into P1 , . . . , Pis , and into Q as well. Thus, on the right-hand side of (4.7) we
multiply numerator and denominator by
s ′
Y
j p, (4.9)
j=1
gcd( j,e
n )=1
where
i ! α−1 β
s′ = p q w−1 ,
M
which comes from the obvious division
n−1 p−1
i = s′ p + i , (4.10)
M M
where s ′ and p−1
M
are integers, by hypothesis. To count the number of elements in the
product (4.9), we do yet another obvious division, namely
! α−1 β ! !
p q w − 1 = p α−1 − 1 q β w + q β w − 1 ,
giving
i ! α−1 i
s′ = p −1 e n − 1).
n + (e (4.11)
M M
Counting the number of elements in the product (4.9) for each of the intervals of length
e
n is no problem; there are exactly φ(en ) elements in each of these intervals. The only
problem is to deal with the remainder term in (4.11), and for that we need Lemma 1.
Using (4.11) and (3.2) with e
n in place of n, we see that the number of elements in the
product (4.9) is
i ! α−1 i i α−1
p − 1 φ(e n) =
n ) + φ(e p φ(e
n ). (4.12)
M M M
But this expression is i A, with A as defined in (4.3). We therefore get from (4.7),
! n−1 P1 · · · Pis · Q
i M n! ≡ Qs ′ (mod e
n ). (4.13)
p iA j
j=1,gcd( j,e
n )=1
Here the bars over the P j and Q indicate that the products (4.8) are taken over all k
relatively prime to e
n , that is,
e
n −1 ie
n −1
Y ! Y !
M

P j := ( j − 1)e
n+k , Q := ise
n+k .
k=1 k=1
gcd(k,e
n )=1 gcd(k,e
n )=1
But then the Gauss-Wilson theorem gives

(
−1 (mod en) if w = 1,
P1 ≡ · · · ≡ Pis ≡ (4.14)
1 (mod e
n) if w > 1.

From the definition of Q we get
ie
n −1 ie
n −1
Y M YM
Q≡ n + k) ≡
(ise k (mod e
n ). (4.15)
k=1 k=1
gcd(k,e
n )=1 gcd(k,e
n )=1
α−1
The Gauss factorial in the denominator of (4.13) can be split up into i p M −1 products
that are congruent to the P j and a remainder that is congruent to Q (mod e n ); this
follows from (4.11). Hence (4.14) and (4.15) together with (4.13) give
! n−1 εB
i M n ! ≡ i A (mod e
n ), (4.16)
p
with A defined by (4.12) and
p α−1 − 1 pα − 1 p α−1 − 1 p−1

B = is − i =i −i = i p α−1 .
M M M M
p−1
Since p is odd, we have ε B = εi M ; this completes the proof of Lemma 2.
We conclude this section with the remark that Theorem 2 is best possible. Indeed,
consider the example M = 3 and n = 70 = 2 · 5 · 7. Here 7 is the only prime factor of
70 that is congruent to 1 (mod 3), and Table 2 shows that 5(3) (3)
1 6 ≡ 52 (mod n). This
is similar to the observation, at the beginning of Section 3, that Lehmer’s result is best
possible.
On the other hand, while Theorem 2 gives a sufficient condition, this condition is
not necessary, as we already saw in Table 4.
5. SOME CONSEQUENCES. We begin this brief section with an immediate con-

sequence of Theorem 2. Let M and n be as in Theorem 2. Since the product of all
partial products 5(M)
j for a fixed M is the Gauss factorial (n − 1)n !, the congruences
(4.1) and the Gauss-Wilson theorem (2.1) give
! n−1 M
M n
! ≡ 1 (mod n).
This implies:
Corollary 1. Let M ≥ 2 be an integer, and suppose that the positive integer n has
at least two distinct prime factors congruent to 1 (mod M). Then the multiplicative
order of the Gauss factorial ( n−1
M n
) ! modulo n is a divisor of M.
While in Lemma 2 the factor w plays only an auxiliary role, collecting all the irrel-
evant prime powers in n (if any), it turns out that we obtain some interesting results if
we consider the Gauss factorial in (4.3) modulo w. Indeed, using the multiplicativity
of φ(n) and the fact that M divides q − 1, we can rewrite A in (4.3) as
p α−1 (q − 1)q β−1

A= φ(q β )φ(w) = p α−1 φ(w) = Cφ(w)
M M
for some integer C. Then, since p ∤ w, we can apply Euler’s generalization of Fermat’s
little theorem and obtain
! φ(w)
pi A = piC ≡ 1 (mod w).

If w > 1 then ε = 1, and the numerator in the congruence in (4.3) is 1. We therefore
get the following corollary.
Corollary 2. Let M, n, and w be as in Lemma 2. Then for i = 1, 2, . . . , M we have

! n−1
i M n ! ≡ 1 (mod w) and 5i(M) ≡ 1 (mod w). (5.1)
This holds for w = 1 because in that case the congruences are trivially true. As a
further consequence we obtain the following result, a special case of which was already
proven as Proposition 4 in [7]. We formulate it as a theorem since it supplements
Theorem 2.
Theorem 3. Let M ≥ 2 be an integer, and suppose that the positive integer n has at
least three distinct prime factors congruent to 1 (mod M). Then
5i(M) ≡ 1 (mod n) for i = 1, 2, . . . , M. (5.2)

α α α
To prove this result, we write n = p1 1 p2 2 p3 3 w, where p1 , p2 , p3 are distinct primes
with p j ≡ 1 (mod M), α j ≥ 1 for j = 1, 2, 3, and gcd( p1 p2 p3 , w) = 1. Now we
α
apply Corollary 2 with w replaced by w j := p j j w, j = 1, 2, 3, obtaining 5i(M) ≡ 1
(mod w j ) for j = 1, 2, 3. The congruences (5.2) then follow immediately from the
Chinese remainder theorem.
In [7] it was shown by similar methods that ( n−1 ) ! = 51(M) ≡ 1 (mod n) under
M n
the conditions of Theorem 3. It was also shown by way of an example that this result,
and thus Theorem 3, is best possible: Let M = 3 and n = 22 · 7 · 13 = 364. Then
obviously 364 ≡ 7 ≡ 13 ≡ 1 (mod 3), and it is easy to compute ( n−1 3 n
) ! = 121364 ! ≡
113 (mod 364).
In analogy to the example at the end of the previous section we show that the
condition in Theorem 3 is not necessary. Let M = 12 and consider n = 146965 =
5 · 7 · 13 · 17 · 19. Then only one prime factor of n is congruent to 1 modulo M, but
5(12)
j ≡ 1 (mod n) for all j.
6. THE GAUSS AND JACOBI BINOMIAL COEFFICIENT THEOREMS. In

Section 1 we remarked in connection with Table 1 that there are no obvious relation-
ships between 5(4) (4)
1 and 52 . One way to explore this further is to consider the quotient
of these partial products. Now, by (1.6) it is clear that for p ≡ 1 (mod 4), 5(4) p−1
1 = 4 !
and 5(4) (4) p−1
1 52 = 2 !, and therefore
p−1 p−1
52(4) 5(4)
1 52
(4)
2
! 2
Q 4 ( p) := = 2 = ! p−1 2 = p−1
. (6.1)
51(4) 5(4) ! 4
1 4
Table 5 lists the values of Q 4 ( p) (mod p) for all the appropriate primes p < 100,
where both the least positive and the least absolute residues are given.
While the least positive residues do not perhaps reveal much, we see a strong con-
nection between the least absolute residues and the representation of p as a sum of
two squares, the existence and uniqueness of which are guaranteed by the well-known
two-squares theorem of Fermat. Even the sign pattern is now quite obvious: the least
absolute residue is positive if and only if a ≡ 1 (mod 4).

Table 5. Q 4 ( p) (mod p) for p ≡ 1 (mod 4), p < 100, and a, b for
which p = a 2 + b2 , with a odd.
p least pos. res. least abs. res. a b
5 2 2 1 2
13 7 −6 3 2
17 2 2 1 4
29 10 10 5 2
37 2 2 1 6
41 10 10 5 4
53 39 −14 7 2
61 10 10 5 6
73 67 −6 3 8
89 10 10 5 8
97 18 18 9 4
All this is in fact explained by the following celebrated theorem of Gauss. We fix
p, a, and b such that
p ≡ 1 (mod 4), p = a 2 + b2 , a ≡ 1 (mod 4). (6.2)
Gauss’s binomial coefficient theorem of 1828 can now be stated as follows.
Theorem 4 (Gauss). Let the prime p and the integer a be as in (6.2). Then
p−1
2
p−1
≡ 2a (mod p). (6.3)
4
As a first easy application of this theorem we show that
5(4) (4)
2 6 ≡ ±51 (mod p) for all p ≡ 1 (mod 4).
Indeed, if this were not the case, then Q 4 ( p) ≡ ±1 (mod p). But by (6.1) and (6.3)
we have Q 4 ( p) ≡ 2a (mod p), so that 2a ≡ ±1 (mod p). The smallest possible solu-
√
tions of this congruence are a = ± p−1
2
. However, by (6.2) we have |a| < p, but also
√
p < p−12
for p > 5. This means that there are no solutions, which was to be shown.
(The case p = 5 is clear from Table 1.)
The analogous investigation for the case M = 3 is a bit more involved, which is
why we present it second. For primes p ≡ 1 (mod 6) we begin, in analogy to (6.1), by
considering
! p−1 p−1
5(3)
2 5(3)
1 52
(3)
2 3 ! 2 3
Q 3 ( p) := (3) = 2 = ! p−1 2 = p−1 . (6.4)
51 5(3) ! 3
1 3
In an attempt to find a congruence for Q 3 ( p) that is analogous to (6.3) one might want
to try another two-squares formula of Fermat, namely p = a 2 + 3b2 for primes p ≡ 1
(mod 6), which is unique up to signs. (For the early history of such representations,
see [15, pp. 14ff.].) However, as columns 2 and 3 of Table 6 show, there seems to be
no apparent relationship modulo p between Q 3 ( p) and a or b.
It was Jacobi who, in 1837, used instead the representation 4 p = x 2 + 3y 2 , which
always has three distinct solutions, namely (|2a|, |2b|), (|a + 3b|, |a − b|), and (|a −

Table 6. Q 3 ( p) (mod p) for p ≡ 1 (mod 6), p < 100, and the solutions of p = a 2 + 3b2 ,
4 p = x 2 + 3y 2 .
p Q 3 ( p) (mod p) a, b x1 , y1 x2 , y2 x1 , y1 r, s
7 −1 2, 1 5, 1 4, 2 1, 3 1, 3
13 5 1, 2 7, 1 5, 3 2, 4 −5,3
19 −7 4, 1 8, 2 7, 3 1, 5 7, 3
31 −4 2, 3 11, 1 7, 5 4, 6 4, 6
37 11 5, 2 11, 3 10, 4 1, 7 −11, 3
43 8 4, 3 13, 1 8, 6 5, 7 −8, 6
61 −1 7, 2 14, 4 13, 5 1, 9 1, 9
67 5 8, 1 16, 2 11, 7 5, 9 −5, 9
73 −7 5, 4 17, 1 10, 8 7, 9 7, 9
79 17 2, 5 17, 3 13, 7 4, 10 −17, 3
97 −19 7, 4 19, 3 14, 8 5, 11 19, 3
3b|, |a + b|). These three solutions, for p < 100, are listed in columns 4–6 of Table 6.
One of these solutions always satisfies y ≡ 0 (mod 3); it is then the corresponding
x, with its sign appropriately chosen, that gives the desired congruence. To be exact,
suppose that the prime p and integers r , s are such that
p ≡ 1 (mod 6), 4 p = r 2 + 3s 2 , r ≡ 1 (mod 3), s ≡ 0 (mod 3). (6.5)
The integer r is then uniquely determined, and we can now state Jacobi’s binomial
coefficient theorem, which is illustrated in the last column of Table 6.
Theorem 5 (Jacobi). Let p and r be as in (6.5). Then

2( p−1)
3
p−1
≡ −r (mod p). (6.6)
3
Proofs of the theorems of Gauss and Jacobi are nonelementary and can be found
in the book [2] by Berndt, Evans, and Williams, which is the standard reference in
the field. For remarks and references, see [2, p. 291]. It is worth giving an explicit
connection between the a in p = a 2 + 3b2 , with its sign fixed by the condition a ≡ −1
(mod 3), and the r as fixed in (6.5):


2a if b ≡ 0 (mod 3),
r = −(a − 3b) if b ≡ 1 (mod 3),

−(a + 3b) if b ≡ 2 (mod 3).
This is an easily obtained modification of congruences in [2, p. 269].

Returning to our observations in Section 1, we now use Jacobi’s theorem to show
that
5(3) (3)
2 6 ≡ 51 (mod p) for all p ≡ 1 (mod 6).
Indeed, if this were not the case, we would have r ≡ −1 (mod p) by (6.4) and (6.6).
Now r = −1 is impossible since r ≡ 1 (mod 3). The next smallest solution, r = p −
√ √
1, is also impossible since by (6.5) we have |r | < 2 p, but we already saw that 2 p <
p − 1 for p > 5.

Now we turn to the congruence 5(3) (3)
2 ≡ −51 (mod p) which, as we saw in Table 1,
does have solutions. Again by (6.4) and (6.6), the congruence is equivalent to r ≡ 1
(mod p). This time we have the solution r = 1, but by the same size argument as
above, there are no others, and (6.5) gives 4 p = 1 + 3s 2 . Now s ≡ 0 (mod 3) and it
also has to be odd, which means that s = 6x + 3 for some positive integer x. If we
substitute this into the expression for 4 p, we get the following result.
Corollary 3. For a prime p ≡ 1 (mod 6) we have 5(3) (3)

2 ≡ −51 (mod p) if and only
if p = 27x 2 + 27x + 7 for an integer x.
The first primes generated by this formula are 7, 61 (see Table 1), 331, 547, and
1951. As is easily seen, negative x give rise to the same primes. It is, of course, not
known whether there are infinitely many primes of this form.
We continue this section with some remarks on congruences for the factorials p−1 4
!
p−1
and 3 !, all of which follow from the theorems of Gauss and Jacobi, respectively.
First, consider primes of the form p ≡ 1 (mod 4). Then
p−1
(a) 4
! ≡ 1 (mod p) only if p = 5.
! p−1 k
(b) 4
! 6≡ −1 (mod p) for k = 1, 2, 4.
! p−1 8
(c) 4
! ≡ −1 (mod p) does hold for p = 17, 241, 3361, 46817, 652081, . . .
For primes of the form p ≡ 1 (mod 6) we have
p−1
(d) 3
! ≡ 1 (mod p) holds for p = 3571, 4219, 13669, 25117, 55897, . . .
! p−1 3
(e) 3
! ≡ 1 (mod p) holds if and only if p = 27x 2 + 27x + 7 for some integer
x.
! p−1 9
(f) 3
! ≡ 1 (mod p) holds if and only if p = 3y 2 + 3y + 1 for some integer
y. Furthermore, the multiplicative order of p−1
3
! (mod p) is 9 if and only if p
2 2
is of the form p = 27x + 9x + 1 or p = 27x + 45x + 19.
! p−1 k
(g) 3
! 6≡ −1 (mod p) for k = 1, 3, 9.
! p−1 18
(h) 3
! ≡ −1 (mod p) holds if and only if p satisfies p 2 = 3y 2 + 3y + 1 for
some integer y. The first few such primes are 13, 181, 2251, 489061.
Statements (c) and (h) are actually connected in the following surprising way: The
identity for p 2 in (h) can be rewritten in the form of the Pell equation (2 p)2 − 3(2y +
1)2 = 1. The infinitely many solutions (An , Bn ) of the equation A2 − 3B 2 = 1 are
given by the recurrence relations (see, e.g., [25, p. 354])
An+2 = 4An+1 − An , A0 = 1, A1 = 2,
Bn+2 = 4Bn+1 − Bn , B0 = 0, B1 = 1.
Then, as is shown in [9], the primes p in (h) are given by primes 12 A2k−1 , while those
2
in (c) are given by prime values of Bn−1 + Bn2 .
While details concerning (h) can be found in [9], statements (a)–(g) are derived and
further discussed in a forthcoming paper [10]. However, some of them follow imme-
diately from results earlier in this section. For instance, if we square (6.1) and use the
fact that by (1.3) we have ( p−1
2
!)2 ≡ −1 (mod p), then we get ( p−1 4
!)4 ≡ −Q 4 ( p)−2
(mod p). Since we know that Q 4 ( p) 6≡ −1 (mod p), this proves statement (b) for
k = 4. The solutions in statement (c) are related to a certain Pell equation; see [4,

p. 318]. These primes form a subsequence of the sequence of all integers a with
the property that a triangle with integer sides (a, a, a − 1) has integer area; see [28,
A103772].
For primes p ≡ 1 (mod 6) we use (1.5) and (1.1) to rewrite (6.4) as
5(3) (3) (3)

1 52 53
! p−1 −3
Q 3 ( p) = ≡− 3
! (mod p). (6.7)
(5(3)
1 )
3
We have seen that Q 3 ( p) ≡ 1 (mod p) is impossible, which proves statement (g) for
k = 3 and also for k = 1. On the other hand, the case Q 3 ( p) ≡ −1 (mod p) is equiv-
alent to Corollary 3; hence (6.7) gives statement (e). A proof along these lines was
suggested by Andrew Granville (private communication with the first author, Decem-
ber, 2004); see also [6].
Statement (e) means that the multiplicative order of p−1 3
! modulo p is 1 or 3 when
2
p = 27x + 27x + 7 for some integer x. Order 1 does actually occur, as statement
(d) indicates. Computations that were kindly carried out for us by Yves Gallot show
that there are 364 such primes up to 109 , while for 762 primes up to 109 the order of
p−1
3
! is 3. It appears to be a difficult question to find a criterion for when the order is
1 and when it is 3. Also, the data suggest that the split between these two classes may
approach 1:2.
Returning to statement (e) and Corollary 3, recall that the polynomial expression
for p comes from
4 p = 1 + 3s 2 = 1 + 3(2y + 1)2 , or p = 3y 2 + 3y + 1 = (y + 1)3 − y 3 , (6.8)
where we have put s = 2y + 1 since s has to be odd. Now Jacobi’s theorem required
3 | s, that is, y = 3x + 1, which led to p = 27x 2 + 27x + 7. In the other two cases,
namely y = 3x and y = 3x + 2, we get p = 27x 2 + 9x + 1 and p = 27x 2 + 45x +
19, respectively. In both these cases the order of p−1 3
! is 9, as is shown in [10]. This,
together with (6.8), gives statement (f).
To conclude this section we note that for Gauss factorials with composite moduli
the situation related to statements (a) and (d) is very different: As Theorem 3 shows,
for each M ≥ 2 we have ( n−1 ) ! ≡ 1 (mod n) for infinitely many n, namely all those
M n
with at least three distinct prime factors p ≡ 1 (mod M).
However, if these composite moduli are prime powers, then the situation remains
very interesting. In fact, in [9] we showed that for a given M ≥ 2 and p ≡ 1 (mod M)
the sequence of multiplicative orders mod p α of the Gauss factorials

p α −1
M
!, α = 1, 2, . . . ,
p
almost always depends in a predictable way on the order of p−1

M
! modulo p. However,
there are some exceptional primes, depending on M, which leads to further interesting
phenomena, mostly explained in [9]. In the case M = 3 this is in fact related to our
remark following statement (h) above.
7. EXTENSIONS OF THE GAUSS BINOMIAL COEFFICIENT THEOREM.

In most of the first five sections of this paper we have dealt with Gauss factorials
and the related products 5(M)
j for composite moduli, often with at least two distinct
prime factors. The case of moduli with only one prime factor, namely prime powers,

turns out to be the most interesting and deepest case in spite of its apparent simplicity.
An indication of this was already provided by the theorems of Gauss and Jacobi in
the previous section, and in the remark at the very end of that section. In the present
section we will describe a further instance of the depth of the prime power case.
In attempting to extend or generalize the theorems of Gauss and Jacobi one might
take two different approaches: First, a natural
! question
is to ask
!2( about congruences
modulo p 2 for the binomial coefficients (( p−1)/2
p−1)/4
and p−1)/3
( p−1)/3
, thus extending the
modulo p congruences of the classical theorems. This was indeed done, as we will see
shortly.
A second approach, natural from the point of view of Gauss factorials, is to consider
the relevant binomial coefficients in terms of factorials. Then one can extend these
objects to Gauss factorials with composite moduli n, and consider the corresponding
quotients modulo n. It turns out that the most interesting case is that of prime power
moduli, which is then related to the first approach, but from a different point of view.
Returning to this first approach, the following extension of the theorem of Gauss to
modulus p 2 was first conjectured by Beukers [3], and later proved by Chowla, Dwork,
and Evans [5].
Theorem 6 (Chowla, Dwork, Evans). Let p and a be as in (6.2). Then

p−1
2 1 p
p−1
≡ 1 + pq p (2) 2a − (mod p 2 ), (7.1)
4
2 2a
where q p (2) := (2 p−1 − 1)/ p is the Fermat quotient to base 2.
Congruences such as (7.1) have been very useful in large-scale computations to

search for Wilson primes, that is, primes p satisfying the congruence ( p − 1)! ≡ −1
(mod p 2 ); see [11] or [12]. For a proof of (7.1) and generalizations to numerous other
binomial coefficients, see [2]. !
Turning now to the Gauss factorial approach, the analogue of (( p−1)/2p−1)/4
(mod p),
α
with modulus n = p , is

p α −1
2
!
(α) p α
B ( p) := ! 2 (mod p ). (7.2)
p α −1
4 p
!
Obviously α = 1 gives the usual binomial coefficient. For α = 2 we were able to show
that the congruence (7.1) is equivalent to
p
B (2) ( p) ≡ 2a − (mod p 2 ). (7.3)
2a
Seeing the simplicity of (7.3) as compared with (7.1), one is led to search numerically
for congruences modulo higher powers of p. Indeed, one readily conjectures that
p p2
B (3) ( p) ≡ 2a − − 3 (mod p 3 ). (7.4)
2a 8a
Continuing with higher powers, we were then able to conjecture and ultimately prove:

Theorem 7. Let p and a be as in (6.2) and let α ≥ 2 be an integer. Then
α
p −1
!
2 p p p2 p α−1
! α 2 ≡ 2a − C0 − C1 3 − · · · − Cα−2
p −1 2a 8a (2a)2α−3
4 p
!

1 2 j − 2 p j
α−1
X
= 2a − 2a 2
(mod p α ), (7.5)
j=1
j j −1 4a
1 2n
!
where Cn := n+1 n
is the nth Catalan number, which is always an integer.
The first few Catalan numbers C0 , C1 , . . . are 1, 1, 2, 5, 14, 42, 132, . . . The proof
of Theorem 7 uses methods similar to those in the proof of Theorem 6; see [2, The-
orem 9.4.3]. The Catalan numbers enter through certain combinatorial identities that
are related to αth powers of particular binomial expressions in complex numbers. For
details, see [8].
If the summation on the right of (7.5) is considered 0 for α = 1, then Gauss’s The-
orem 4 can be seen as a special case of (7.5). We already remarked that for α = 2
the congruences (7.5) and (7.1) are equivalent. This leads to the natural question of
whether one can derive a binomial coefficient analogue to (7.5) for α = 3. This is in
fact possible, and we obtain the following mod p 3 extension of Theorem 6.
Theorem 8. Let p and a be as in (6.2). Then

p−1
2 p p2
p−1
≡ 2a − − 3
4
2a 8a
! !
× 1 + 12 pq p (2) + 18 p 2 2E p−3 − q p (2)2 (mod p 3 ), (7.6)
where E n is the nth Euler number.
For further details and proofs, see [8], where Jacobi’s Theorem 5 has also been ex-
tended in a similar fashion.
To summarize: Comparing the congruences (7.5) and (7.6), it is clear that for higher
congruences the Gauss factorials with prime power moduli are the more natural objects
to study than the usual factorials.
8. CONCLUSION. The number-theoretic object we propose to call a Gauss facto-

rial is certainly not new. For instance, it has played an important role in the study of
arithmetic properties of binomial coefficients [18], and is essential in the definition of
Morita’s p-adic gamma function (see, e.g., [2, p. 277]). This paper, however, is a study
of Gauss factorials as objects in their own right. In this study, we relied heavily on
numerical experimentation using the computer algebra system Maple. In fact, without
the assistance of such a tool this study would not have been possible.
Our purpose in this paper has been threefold: First, to study the special Gauss fac-
torials ( n−1
M n
) !, and in particular their values and multiplicative orders modulo n, in the
spirit of the remarkable but little-known Gauss-Wilson theorem, which is the special
case M = 1.
The second purpose has been the introduction and study of the associated partial
products 5(M) j , defined by (2.5), which extend the Gauss factorials since 51
(M)
=

( n−1
M n
) !. Our main results, Theorems 2 and 3, are different in nature from most re-
sults in classical number theory in that they depend on the number of different prime
factors of a given n. The only result of this nature of which we are aware is that of
D. H. Lehmer (Lemma 1); our results complement his in that they concern the prod-
ucts of “totatives,” as opposed to their numbers.
Our third purpose has been to show that some deep extensions of the binomial
theorems of Gauss and Jacobi appear in a simpler and more natural way, and can be
further extended, if stated in terms of Gauss factorials. This also points to the fact that
Gauss factorials are particularly worthwhile objects to study when the modulus n is a
power of an odd prime.
In summary, we hope that we have demonstrated the inherent beauty, depth, and
usefulness of Gauss factorials.
ACKNOWLEDGMENTS. We would like to thank Yves Gallot for generously carrying out the computations
reported in Section 6. We also thank Andrew Granville for suggesting the proof mentioned in Section 6.
REFERENCES
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861. http://dx.doi.org/10.1090/S0025-5718-98-00951-X
2. B. C. Berndt, R. J. Evans, and K. S. Williams, Gauss and Jacobi Sums, Wiley, New York, 1998.
3. F. Beukers, Arithmetical properties of Picard-Fuchs equations, in Seminar on Number Theory—Paris,
1982–83, Progr. Math., vol. 51, M.-J. Bertin and C. Goldstein, eds., Birkhäuser, Boston, MA, 1984,
33–38.
4. J. M. Borwein and D. Bailey, Mathematics by Experiment. Plausible Reasoning in the 21st Century, 2nd
ed., A K Peters, Wellesley, MA, 2008. !
5. S. Chowla, B. Dwork, and R. Evans, On the mod p 2 determination of (( p−1)/2 p−1)/4 , J. Number Theory 24
(1986) 188–196. http://dx.doi.org/10.1016/0022-314X(86)90102-2
6. J. B. Cosgrave, Trinity College Dublin Mathematical Society Lecture, available at http://staff.spd.
dcu.ie/johnbcos/jacobi.htm.
7. J. B. Cosgrave and K. Dilcher, Extensions of the Gauss-Wilson theorem, Integers 8 (2008) A39; available
at http://www.integers-ejcnt.org/vol8.html.
8. , Mod p 3 analogues of theorems of Gauss and Jacobi on binomial coefficients, Acta Arith. 142
(2010) 103–118. http://dx.doi.org/10.4064/aa142-2-1
9. , The multiplicative orders of certain Gauss factorials, Int. J. Number Theory 7 (2011) 145–171.
http://dx.doi.org/10.1142/S179304211100396X
10. , The Gauss-Wilson theorem for one-third, one-quarter and one-sixth intervals (in preparation).
11. R. Crandall, Topics in Advanced Scientific Computation, Springer-Verlag, New York, 1996.
12. R. E. Crandall, K. Dilcher, and C. Pomerance, A search for Wieferich and Wilson primes, Math. Comp.
66 (1997) 433–449. http://dx.doi.org/10.1090/S0025-5718-97-00791-6
13. L. E. Dickson, History of the Theory of Numbers. Volume I: Divisibility and Primality, Chelsea, New
York, 1966.
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Verlag, New York, 1977.
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org/10.4153/CMB-1958-008-7
17. C. F. Gauss, Disquisitiones Arithmeticae (trans. and preface by A. A. Clarke), Yale University Press,
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W. C. Waterhouse, Springer-Verlag, New York, 1986.
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ers, in Organic Mathematics—Burnaby, BC, 1995, CMS Conf. Proc., vol. 20, American Mathematical
Society, Providence, RI, 1997, 253–276.
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//dx.doi.org/10.4153/CMB-2002-012-1

20. G. H. Hardy and E. M. Wright, An Introduction to the Theory of Numbers, 5th ed., Oxford University
Press, New York, 1979.
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10.4153/CJM-1955-038-5
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http://dx.doi.org/10.2307/2312481
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Wiley, New York, 1991.
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395.
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JOHN B. COSGRAVE was born in Bailieboro, County Cavan, Ireland. He received his B.Sc. (1968) and
Ph.D. (1972) in Mathematics from Royal Holloway College (London University). He worked at RHC,
Manchester, Jos (Nigeria), Carysfort College (Dublin), and finally St. Patrick’s College, Drumcondra, Dublin,
where—shortly before retiring in 2007—he had the pleasure of having Doron Zeilberger as his department’s
international external assessor. Besides elementary number theory, his interests include reading (all kinds),
music, and cycling, and—together with his wife Mary, whom he met at RHC—he is a daily swimmer in
Dublin Bay.
79 Rowanbyrn, Blackrock, County Dublin, Ireland
jbcosgrave@gmail.com
http://staff.spd.dcu.ie/johnbcos/
KARL DILCHER received his undergraduate education at the Technische Universität Clausthal in Germany.
He then did his graduate studies at Queen’s University in Kingston, Ontario, and finished his Ph.D. there
in 1983 under the supervision of Paulo Ribenboim. He is currently a professor at Dalhousie University in
Halifax, Nova Scotia, Canada, where he first arrived in 1984 as a postdoctoral fellow. His research interests
include classical analysis, special functions, and elementary and computational number theory.
Department of Mathematics and Statistics, Dalhousie University, Halifax, Nova Scotia, B3H 4R2, Canada
dilcher@mathstat.dal.ca

NOTES
Edited by Ed Scheinerman
Enumerating the Rationals from

Left to Right
S. P. Glasby
Abstract. The rationals can be enumerated using Stern-Brocot sequences SBn , Calkin-Wilf
sequences CW n , or Farey sequences Fn . We show that the rationals in SBn , CW n , and Fn can
be computed using very similar second-order linear recurrence relations. (This, incidentally,
obviates the need to compute previous sequences SBi , CW i , Fi with i < n.)
1. INTRODUCTION. There are three well-known sequences used to enumerate the

rationals: the Stern-Brocot sequences SB n , the Calkin-Wilf sequences CW n , and the
Farey sequences Fn . The purpose of this note is to show that all three sequences can
be constructed (left-to-right) using almost identical recurrence relations. The Stern-
Brocot (S-B) and Calkin-Wilf (C-W) sequences give rise to complete binary trees
related to Figure 1. These trees have many beautiful algebraic, combinatorial, com-
putational, and geometric properties [2, 5, 4]. Well-written introductions to the S-B
tree and Farey sequences can be found in [3], and to the C-W tree in [2]. We shall
focus on sequences rather than trees.
a c p
b d q
a +c p p+q
b+d p+q q
Figure 1. S-B rules (left), and C-W rules (right).
Two fractions ab < dc are called adjacent if bc − ad = 1. Adjacent fractions are

necessarily reduced, i.e., gcd(a, b) = gcd(c, d) = 1. The mediant of ab < dc is b+d
a+c
.
a c a a+c c
A short calculation shows that if b < d are adjacent, then b < b+d < d are pairwise
0 1 (and hence reduced). The sequences SB n are defined recursively: SB 0 =
adjacent
, represents 0 and ∞ as reduced fractions, and SB n is computed from SB n−1 by
1 0
inserting mediants between consecutive fractions. Thus

0 1 1
SB 1 = , , ,
1 1 0

0 1 1 2 1
SB 2 = , , , , ,
1 2 1 1 0

0 1 1 2 1 3 2 3 1
SB 3 = , , , , , , , , , ....
1 3 2 3 1 2 1 1 0

A simple induction shows that |SB n | = 2n + 1. Thus 2n−1 mediants are inserted into
SB n−1 to form SB n . The C-W sequences are defined using the right rule in Figure 1:

1
CW 1 : = ,
1

1 2
CW 2 = , ,
2 1

1 3 2 3
CW 3 = , , , ,
3 2 3 1

1 4 3 5 2 5 3 4
CW 4 = , , , , , , , ,....
4 3 5 2 5 3 4 1
A simple induction shows that |CW n | = 2n−1 . Another simple induction (see [2,
p. 360]) shows that the fractions in CW n have the form

b−1 b0 b N −2
CW n = , ,..., where N = 2n−1 ,
b0 b1 b N −1
and the denominator of a given fraction is the numerator of the succeeding fraction.
Indeed, this property obtains even when the sequences CW 1 , CW 2 , CW 3 , . . . are con-
catenated to form
z}|{ z }| { z }| {
1 1 2 1 3 2 3 1 4 c0 c1 c2
CW ∞ := , , , , , , , , . . . , , . . . = , , ,... .
1 2 1 3 2 3 1 4 1 c1 c2 c3
p
The Farey sequence of order n contains all the reduced fractions q
with 0 ≤ p ≤
q ≤ n, in their natural order. Thus

0 1
F1 = , ,
1 1

0 1 1
F2 = , , ,
1 2 1

0 1 1 2 1
F3 = , , , , ,
1 3 2 3 1

0 1 1 1 2 3 1
F4 = , , , , , , ,....
1 4 3 2 3 4 1
A standard way to compute Fn from Fn−1 is to insert mediants between consecutive
fractions of Fn−1 only when this gives a denominator of size n (see [3, p. 118]). Thus
Fn is a subsequence of SB n . The number of reduced fractions an with 1 ≤ a < n is
denoted ϕ(n). Counting the P nonzero reduced fractions with denominator j, together
with zero, gives |Fn | = 1 + nj=1 ϕ( j). The mediant rule above implies that consecu-
tive fractions in SB n and Fn are adjacent (see also [3, p. 119]).
S
2. RECURSIVE RESULTS. It is shown in [3] and [2] that SB ∞ = ∞ n=0 SB n and
CW ∞ contain every (reduced) positive rational precisely once. Although SB n , CW n ,
and Fn are defined “top-down” they can be computed from “left to right” via almost
identical recurrence relations. The 2-exponent ν2 (i) of a positive integer i is defined
by i = 2ν2 (i) j, where j is an odd integer.
November 2011] NOTES 831

h i
a−1 a0 a1 a
Theorem 1. Write SB n = , , , . . . , b NN −1
b−1 b0 b1 −1
where N = 2n . Then
a−1 = 0, a0 = 1, ai = ki ai−1 − ai−2 for 1 ≤ i < N , (1a)

b−1 = 1, b0 = n, bi = ki bi−1 − bi−2 for 1 ≤ i < N , (1b)
where ki = 2ν2 (i) + 1.
Theorem 2. Write

a0 a1 ai−1
CW ∞ = , ,..., ,...
a1 a2 ai
and

b−1 b0 b1 b N −2
CW n = , , ,...,
b0 b1 b2 b N −1
where N = 2n−1 . Then the ai and bi can be computed via the recurrence relations
a−1 = 0, a0 = 1, ai = ki ai−1 − ai−2 for 1 ≤ i < ∞, (2a)

b−1 = 1, b0 = n, bi = ki bi−1 − bi−2 for 1 ≤ i < N , (2b)
where ki = 2ν2 (i) + 1.

h i
A−1 A0 A1 A −1
Theorem 3. Write the Farey sequence Fn of order n as Fn = B−1 , B0 , B1 , . . . , B NN −1 .
Then the numerators Ai and the denominators Bi can be computed via the recurrence
relations
A−1 = 0, A0 = 1, Ai = K i Ai−1 − Ai−2 for 1 ≤ i < N , (3a)

B−1 = 1, B0 = n, Bi = K i Bi−1 − Bi−2 for 1 ≤ i < N , (3b)
j k P
Bi−2 +n
where K i = Bi−1
, and N = nj=1 ϕ( j).
To illustrate Theorem 1, SB 4 can be computed from left to right using Table 1.

Note that the numbers ki are the same as the numbers ki′ generated by the recurrence
k1′ = 1, k2′ j+1 = 1, k2′ j = k ′j + 2 for j ≥ 0. (Proof by induction: k1 = k1′ and k2 j+1 = 1,
k2 j = k j + 2 hold for j ≥ 1 as ν2 (2 j + 1) = 0 and ν2 (2 j) = ν2 ( j) + 1. Thus ki = ki′
for all i ≥ 1.)
Table 1. Computing SB4 from left to right.
i −1 0 1 2 3 4 5 6 7 8 9 10 11 12 13 14 15
ai 0 1 1 2 1 3 2 3 1 4 3 5 2 5 3 4 1
bi 1 4 3 5 2 5 3 4 1 3 2 3 1 2 1 1 0
ki 1 3 1 5 1 3 1 7 1 3 1 5 1 3 1

Proof of Theorem 1. Our proof uses induction on n. It suffices to prove (1a) as the
proof of (1b) is similar (just change the as to bs). Clearly (1a) is true for n = 0 as
a′ a0′ a ′N /2−1
SB 0 = 01 , 10 . Assume n > 0 and (1a) is true for SB n−1 . Let b−1
′ , b′ , . . . , b′ be the
−1 0 N /2−1
fractions in SB n−1 . The way mediants are inserted to create SB n is shown in Figure 2.
Dotted lines denote the repetition of a fraction, and dashed lines denote the formation
of a mediant. The repetition of fractions means
a2 j−1 = a ′j−1 and b2 j−1 = b′j−1 for 0 ≤ j < N /2, (4)
and the formation of mediants means
a2 j = a ′j−1 + a ′j and b2 j = b′j−1 + b′j for 0 ≤ j < N /2. (5)
Certainly the formulas for a−1 and a0 in (1a) are correct, as SB n starts with 01 , n1 .
Suppose now that 1 ≤ i < N , and consider the cases when i is even and odd.
a ′j−2 a ′j−1 a ′j
SBn−1 b′j−2 b′j−1 b′j
a2 j−3 a2 j −2 a2 j−1 a2 j a2 j+1

SBn b2 j−3 b2 j −2 b2 j−1 b2 j b2 j+1
Figure 2. Constructing SBn from SBn−1 by inserting mediants.
Case 1. i = 2 j is even and j ≥ 1. The following shows that (1a) holds for even i:
k2 j a2 j−1 − a2 j−2 = (k j + 2)a2 j−1 − a2 j−2 as k2 j = k j + 2,

= (k j + 2)a ′j−1 − (a ′j−2 + a ′j−1 ) by (4) and (5),
= k j a ′j−1 − a ′j−2 + a ′j−1 canceling a ′j−1 ,
= a ′j + a ′j−1 as a ′j = k j a ′j−1 − a ′j−2 by
induction,
= a2 j by (5).
Case 2. i = 2 j + 1 is odd and j ≥ 1. It follows from k2 j+1 = 1 and (4) and (5) that
k2 j+1 a2 j − a2 j−1 = a2 j − a2 j−1 = (a ′j−1 + a ′j ) − a ′j−1 = a ′j = a2 j+1 ,
as desired. This completes the induction on n.
A different (and very interesting) method for computing terms of SB n is given

in [1]. It uses continued fraction expansions and “normalized additive factorizations.”
As the recurrence (1a) is independent of n, the numerators for SB n−1 reappear as the
first 2n−1 + 1 numerators for SB n . We now show that (half of) the denominators bi in
SB n reappear (remarkably!) in CW n , and the numerators ai also reappear in CW ∞ .
Accordingly, we have used the same notation ai , bi in Theorem 2 as in Theorem 1.

Proof of Theorem 2. Figure 3 shows that the numbers ai must satisfy the recurrence
relation:
(6.1) (6.2)
a0 = 1, a2 j−1 = a j−1 and a2 j = a j−1 + a j for j > 0. (6)
Define a−1 = 0. With this definition (6.1) and (6.2) also hold when j = 0. We shall
now prove, by induction on i, that the formulas for ai in (2a) are correct. This is clear
for i = −1, 0. Now suppose i > 0, and assume the formulas in (2a) hold for all i ′ < i.
a0 1
a1
= 1
a j−1
aj
a1 1 a2 2
a2
= 2 a3
= 1
a j−1 a j−1 +a j
a j−1 +a j aj
a3 1 a4 3 a5 2 a6 3
a4
= 3 a5
= 2 a6
= 3 a7
= 1
Figure 3. C-W tree (left), and C-W rules (right).
C ASE 1. i = 2 j where j ≥ 1. Then
k2 j a2 j−1 − a2 j−2 = (k j + 2)a2 j−1 − a2 j−2 as k2 j = k j + 2,

= (k j + 2)a j−1 − (a j−2 + a j−1 ) by (6.1) and (6.2),
= k j a j−1 − a j−2 + a j−1 canceling a j−1 ,
= a j + a j−1 by induction on i,
= a2 j by (6.1).
C ASE 2. i = 2 j − 1 where j ≥ 1. Then k2 j+1 = 1, so (6) implies
k2 j+1 a2 j − a2 j−1 = a2 j − a2 j−1 = (a j−1 + a j ) − a j−1 = a j = a2 j+1 .
This completes the inductive proof of (2a).

b
The proof of (2b) is now straightforward. The first term b−1 0
= n1 of CW n equals
a N −1
aN
where N = 2n−1 . As CW n is a consecutive subsequence of CW ∞ , it fol-
lows that a N +i = bi for −1 ≤ i < N . Now a N −1 = 1, a N = n, and (2a) implies
that a N +i = k N +i a N +i−1 − a N +i−2 for 1 ≤ i < N . However, ν2 (N + i) = ν2 (i) and
k N +i = ki for 1 ≤ i < N . Thus (2b) follows by replacing a N +i with bi , and k N +i with
ki , for 1 ≤ i < N .
Theorem 3 is previously known (see Exercise 4.61 in [3, p. 150]). We include

Theorem 3 and its proof both for comparison with Theorems 1 and 2, and for the
reader’s convenience.
Proof of Theorem 3. As the first two fractions of Fn are 01 and n1 , the recurrences
Ai−2 Ai−1
(3a, b) are correct for i = −1, 0. Now suppose 1 ≤ i < N . By definition, Bi−2 , Bi−1 ,
Ai
and Bi
are consecutive terms in Fn . We shall prove that Ai = a and Bi = b, where

a = K i Ai−1 − Ai−2 and b = K i Bi−1 − Bi−2 . We know that Ai−1 Bi−2 − Bi−1 Ai−2 = 1
since consecutive Farey fractions are adjacent. Hence
a Bi−1 − b Ai−1 = (K i Ai−1 − Ai−2 )Bi−1 − (K i Bi−1 − Bi−2 )Ai−1

= Ai−1 Bi−2 − Bi−1 Ai−2 = 1. (7)
B +n B +n
Consider the inequalities i−2
Bi−1
− 1 < K i ≤ i−2 Bi−1
. Multiplying by Bi−1 and sub-
tracting Bi−2 gives n − Bi−1 < b ≤ n. It follows from Bi−1 ≤ n that 0 < b ≤ n, and it
Ai−2 Ai−1
follows from (7) that Bi−2 < Bi−1 < ab . Since ABii is, by definition, the next fraction in
Fn after ABi−1
i−1 Ai−1
, we conclude that Bi−1 < ABii ≤ ab .
We must show that ABii = ab . Suppose to the contrary that Ai
Bi
< ab . Then
(8.1) (8.2)
1 ≤ a Bi − b Ai and 1 = Ai Bi−1 − Bi Ai−1 . (8)
Multiplying (8.1) by Bi−1 , and (8.2) by b, and then adding gives

(7)
n < Bi−1 + b ≤ (a Bi − b Ai )Bi−1 + b(Ai Bi−1 − Bi Ai−1 ) = (a Bi−1 − b Ai−1 )Bi = Bi .
This is a contradiction since ABii ∈ Fn has Bi ≤ n. Hence ABii = ab . As both fractions are
reduced (and Ai , Bi , b > 0), we conclude that Ai = a and Bi = b, as desired.
3. CONCLUDING REMARKS. The On-Line Encyclopedia of Integer Sequences [6]

has a wealth of useful information about the sequences a0 , a1 , a2 , . . . (A002487) and
k1 , k2 , k3 , . . . (A037227). However, the connection in Theorem 2 between these se-
quences is new. Note that an counts the number of ways that n can be written as a
sum of powers of 2, each power being used at most twice. For example, a4 = 3 as
22 = 2 + 2 = 2 + 1 + 1.
ACKNOWLEDGMENTS. The author would like to thank the referees for their insightful comments.
REFERENCES
1. B. Bates, M. Bunder, and K. Tognetti, Locating terms in the Stern-Brocot tree, European J. Combin. 31
(2010) 1020–1033. http://dx.doi.org/10.1016/j.ejc.2007.10.005
2. N. Calkin and H. S. Wilf, Recounting the rationals, Amer. Math. Monthly 107 (2000) 360–363. http:
//dx.doi.org/10.2307/2589182
3. R. L. Graham, D. E. Knuth, and O. Patashnik, Concrete Mathematics, 2nd ed., Addison-Wesley, Reading,
MA, 1994.
4. M. Hockman, Continued fractions and the geometric decomposition of modular transformations, Quaest.
Math. 29 (2006) 427–446. http://dx.doi.org/10.2989/16073600609486174
5. M. Niqui, Exact arithmetic on the Stern-Brocot tree, J. Discrete Algorithms 5 (2007) 356–379. http:
//dx.doi.org/10.1016/j.jda.2005.03.007
6. On-Line Encyclopedia of Integer Sequences, available at http://oeis.org.
Department of Mathematics, Central Washington University, Ellensburg, WA 98926-7424

http://www.cwu.edu/∼glasbys/

On the Banach-Steinhaus Theorem
for Topological Groups
Semyon N. Litvinov
Abstract. We prove the Banach-Steinhaus theorem and some of its consequences under the
assumption that the underlying space is a topological group of second Baire category.
1. PRELIMINARIES. Let us start with a definition.
Definition. Let X and Y be topological groups. A family of homomorphisms aλ :

X → Y , λ ∈ 3, is said to be equicontinuous if, for every open neighborhood W of the
identity element in Y , there is an open neighborhood V of the identity element in X
such that aλ (V ) ⊂ W for all λ ∈ 3.
The classical Banach-Steinhaus theorem may be stated as follows.
Theorem 1. Let X be a Banach space, and let (Y, k · k) be a normed space. If aλ :

X → Y is a continuous linear map for every λ ∈ 3 such that supλ∈3 kaλ (x)k < ∞
for all x ∈ X , then the family {aλ }λ∈3 is equicontinuous.
In [5, p. 44], Theorem 1 is proved under the assumptions that X and Y are topo-
logical vector spaces, and the set of all elements of X with bounded orbits is a set of
second Baire category in X . In the next section we will extend Theorem 1 to the case
where X is a topological group of second Baire category.
Let (X, ∗, τ ) be a topological group. Recall that if (X, τ ) is Hausdorff, then it is of
second Baire category in each of the following cases [5, p. 43]:
(a) X is complete and metrizable;
(b) X is locally compact.
Also, if (X, ∗, τ ) is a Hausdorff topological group, then it is metrizable if and only
if the topological space (X, τ ) has a countable base of neighborhoods of the identity
element [1].
The identity element of a topological group (X, ∗, τ ) will be denoted by e, the
inverse of x ∈ X by x −1 , while x m will stand for |x ∗ ·{z
· · ∗ x}. For basic properties and
examples of topological groups, see [1]. m times
2. BANACH-STEINHAUS THEOREM FOR TOPOLOGICAL GROUPS. The

following is a generalization of the Banach-Steinhaus theorem to the case in which the
underlying space is a topological group.
Theorem 2. Let (X, ∗, τ ) be a topological group of second Baire category, and let
(Y, k · k) be a normed space. If aλ : X → Y is a continuous homomorphism for ev-
ery λ ∈ 3 such that supλ∈3 kaλ (x)k < ∞ for all x ∈ X , then the family {aλ }λ∈3 is
equicontinuous.
836
c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118
Proof. Fix ǫ > 0. For L ∈ N, let B L = {y ∈ Y : kyk ≤ L} and put
\
XL = aλ−1 (B L ).
λ∈3
Then X L is closed for every L. Also, one can easily check that
[
X= XL.
L
Using the assumption that X is of second Baire category, we infer that there exists
L 0 such that X L 0 contains a nonempty open set. Therefore, aλ (U ) ⊂ B L 0 , λ ∈ 3, for
some nonempty open U ⊂ X . This allows us to conclude that
â(x) = sup kaλ (x)k ≤ L 0 for every x ∈ U.

λ∈3
Note that, by the triangle inequality and the fact that each aλ is a homomorphism, we
have
â(x ∗ z) ≤ â(x) + â(z) for all x, z ∈ X.
Pick x0 ∈ U , and let V be an open neighborhood of e such that x0 ∗ V ⊂ U . If x ∈ V ,

then x0 ∗ x ∈ U , which together with â(x0−1 ) = â(x0 ) yields
â(x) ≤ â(x0−1 ) + â(x0 ∗ x) ≤ 2L 0 .
Now let m ∈ N be such that 2Lm0 < ǫ, and let O be an open neighborhood of e for
which O m ⊂ V . If x ∈ O, then x m ∈ V , and we obtain
m · â(x) = â(x m ) ≤ 2L 0 ,
which implies that â(x) < ǫ. Thus, for a given ǫ > 0, we were able to find an open
neighborhood O of e such that kaλ (x)k < ǫ for all λ ∈ 3, that is, the family {aλ }λ∈3
is equicontinuous.
3. APPLICATIONS. The theorem below is, in a sense, a “uniform” prototype of

the Banach principle on almost everywhere convergence of sequences of measurable
functions (see [2, p. 63]; also [3]).
Theorem 3. Let (X, ∗, τ ) be a topological group of second Baire category, and let
(Y, k · k) be a Banach space. For a sequence an : X → Y of continuous homomor-
phisms, consider the conditions:
(A) the sequence {an (x)} converges for every x ∈ X ;
(B) â(x) = supn kan (x)k < ∞ for every x ∈ X ;
(C) the operator â : X → R is continuous at e;
(D) the set C = {x ∈ X : {an (x)} converges} is closed in X .
Then the implications (A)⇒(B)⇒(C)⇒(D) hold. If, in addition, the sequence {an (x)}
converges for every x in a dense subset of X , then conditions (A)–(D) are equivalent.

Proof. The implication (A)⇒(B) is obvious, while the implication (B)⇒(C) is a par-
ticular case of Theorem 2.
(C)⇒(D): Take x̄ ∈ C and fix ǫ > 0. Using continuity of â at e, we can find an
open neighborhood O of e such that
ǫ
â(x) < for all x ∈ O.
3
Next, there is an open neighborhood V of x̄ satisfying the condition x̄ −1 ∗ V ⊂ O. In
this case, it is possible to find z 0 ∈ C for which x0 = x̄ −1 ∗ z 0 ∈ O. So, we have
ǫ
sup kan (x0 )k = â(x0 ) < .
n 3
Further, since x̄ ∗ x0 = z 0 ∈ C, there exists N ∈ N such that

ǫ
kam (x̄ ∗ x0 ) − an (x̄ ∗ x0 )k < whenever m, n ≥ N .
3
Now, for m, n ≥ N , we have
kam (x̄) − an (x̄)k ≤ kam (x̄ ∗ x0 ) − an (x̄ ∗ x0 )k + kam (x0 )k + kan (x0 )k < ǫ,
which implies that the sequence {an (x̄)} is convergent. Thus C = C.

The remaining part of the proof is straightforward.
Theorem 4. Let (Y, k · k) be a normed space, and let X and {an } be as above. If,
for every x ∈ X , there exists yx ∈ Y satisfying limn→∞ kan (x) − yx k = 0, then the
operator a : X → Y given by the formula a(x) = yx , x ∈ X , is a continuous homo-
morphism.
Proof. Because the sequence {an (x)} converges for every x ∈ X , condition (A) of
Theorem 3 is satisfied, which implies that the operator â is continuous at e. Therefore,
since a(x) = limn→∞ an (x) and, consequently,
ka(x)k = lim kan (x)k ≤ sup kan (x)k = â(x),

n→∞ n
the operator a is also continuous at e, hence on X .
ACKNOWLEDGMENTS. The author is grateful to Professor Vyacheslav Chistyakov of the State University
Higher School of Economics, Nizhny Novgorod, for posting a question that is answered in Theorem 4 of the
present article.
REFERENCES
1. E. Hewitt and K. A. Ross, Abstract Harmonic Analysis, vol. I, Springer-Verlag, Berlin, 1963.
2. U. Krengel, Ergodic Theorems, Walter de Gruyter, Berlin, 1985.
3. S. Litvinov, The Banach principle for topological groups, Atti Sem. Mat. Fis. Univ. Modena e Reggio
Emilita LIII (2005) 323–330.
4. H. L. Royden and P. M. Fitzpatrick, Real Analysis, 4th ed., Prentice Hall, Boston, 2010.
5. W. Rudin, Functional Analysis, 2nd ed., McGraw Hill, New York, 1991.
Department of Mathematics, Pennsylvania State University, Hazleton, PA 18202

snl2@psu.edu
838
c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118
Ramanujan’s “Most Beautiful Identity”
Michael D. Hirschhorn
Abstract. We give a simple proof of the identity which for Hardy represented the best of
Ramanujan. On the way, we give a new proof of an important identity that Ramanujan stated
but did not prove.
Of all the 4000 or so identities Ramanujan presented, Hardy chose one which for
him represented the best of Ramanujan. I would like to show you this identity, and
prove it.
Following Euler, we define a partition of the positive integer n as a representation
of n as a sum of positive integers, in which order is unimportant. The partitions of 4
are 4 = 3 + 1 = 2 + 2 = 2 + 1 + 1 = 1 + 1 + 1 + 1. The number of partitions of n
is denoted by p(n); thus, p(4) = 5. For convenience, we define p(0) = 1.
Euler showed that the partition generating function
X
P(q) = p(n)q n = 1 + 1q + 2q 2 + 3q 3 + 5q 4 + · · ·
n≥0
satisfies
1
p(q) = ,
(q; q)∞
where
Y
(a; q)∞ = (1 − aq n ).
n≥0
He also showed that
(q; q)∞ = 1 − q − q 2 + q 5 + q 7 − q 12 − q 15 + + − − · · · .
In the series on the right, the terms occur in pairs, alternately with coefficients −1
and +1. The powers, 1, 2, 5, 7, 12, 15, . . ., are known as the pentagonal numbers, and
Euler’s expansion above as the pentagonal number theorem. The easiest way of gener-
ating the pentagonal numbers is to proceed as follows.
The triangular numbers 1, 1 + 2, 1 + 2 + 3 and so on are
1, 3, 6, 10, 15, 21, 28, 36, 45, 55, 66, 78, . . . .
Two of every three are divisible by 3. If we divide these by 3 we obtain the pentagonal
numbers!
A beautiful combinatorial proof of Euler’s pentagonal number theorem was given
by F. Franklin in 1881, and is reproduced in Hardy and Wright [3].
Euler’s pentagonal number theorem is the special case a = 1 of Jacobi’s triple prod-
uct identity when written

(a −1 q; q 3 )∞ (aq 2 ; q 3 )∞ (q 3 ; q 3 )∞
= 1 − a −1 q − aq 2 + a −2 q 5 + a 2 q 7 − a −3 q 12 − a 3 q 15 + + − − · · · .
A proof of the triple product identity is given in the Appendix. Also, it is important
that I mention that Ramanujan found a marvellous and powerful extension of the triple
product identity.
Back to the main story: We have that the partition generating function is the recip-
rocal of the series 1 − q − q 2 + q 5 + q 7 − − + + · · · . This implies
p(0) = 1,
p(1) − p(0) = 0,
p(2) − p(1) − p(0) = 0,
p(3) − p(2) − p(1) = 0,
and, more generally, for n > 0,
p(n) − p(n − 1) − p(n − 2) + p(n − 5) + p(n − 7) − − + + · · · = 0.
Here, you recognise the pattern of + and − signs, and the numbers 1, 2, 5, 7, . . . . For
each n, the sum on the left terminates, since all terms with negative argument are zero.
P. MacMahon, who was in Cambridge with Hardy and Ramanujan, used the above
recurrence to calculate p(n) for n ≤ 200, and serendipitously listed the values in
groups of five thus:
1 7 42 176 627 1958 · · ·

1 11 56 231 792 2436 · · ·
2 15 77 297 1002 3010 · · ·
3 22 101 385 1255 3718 · · ·
5 30 135 490 1575 4565 · · ·
Ramanujan observed that the numbers at the bottom of each group are divisible by 5,
that is, 5| p(5n + 4). He also observed that 7| p(7n + 5), 11| p(11n + 6), and, on the
basis of the very small amount of evidence provided by MacMahon’s table, formulated
a very general conjecture, which was essentially correct; the proof was completed by
Oliver Atkin in 1967.
Ramanujan did much more than prove 5| p(5n + 4). We can write down the gener-
ating function of the p(5n + 4),
X
p(5n + 4)q n = 5 + 30q + 135q 2 + 490q 3 + 1575q 4 + 4565q 5 + · · · .
n≥0
Ramanujan [4, p. 213] makes the claim that this series can be written as a neat product,
X
n(q 5 ; q 5 )5∞
p(5n + 4)q = 5 . (1)
n≥0
(q; q)6∞
Hardy [3, p. xxxv] says of (1): “It would be difficult to find more beautiful formulae
than the ‘Rogers-Ramanujan’ identities, but here Ramanujan must take second place
to Rogers; and, if I had to select one formula from all Ramanujan’s work, I would

agree with Major MacMahon in selecting [(1)].” Hence, I refer to (1) as “Ramanujan’s
most beautiful identity.”
1
P n
P n (q 5 ;q 5 )5∞
In short (1) says, if (q;q)∞
= p(n)q then p(5n + 4)q = 5 (q;q)6∞
. Note that
this formulation does not require that we have a combinatorial interpretation of the
p(n); the statement can be regarded as purely algebraic.
My aim is to sketch a proof of the wonderful identity (1).
We can write, with ω 6 = 1 a fifth root of unity,
1 (ωq; ωq)∞ (ω2 q; ω2 q)∞ (ω3 q; ω3 q)∞ (ω4 q; ω4 q)∞

= . (2)
(q; q)∞ (q; q)∞ (ωq; ωq)∞ (ω2 q; ω2 q)∞ (ω3 q; ω3 q)∞ (ω4 q; ω4 q)∞
The denominator of the right-hand side of (2) is
Y
(1 − q n )(1 − ωn q n )(1 − ω2n q 2n )(1 − ω3n q 3n )(1 − ω4n q 4n )
n≥1
Y Y
= (1 − q n )5 · (1 − q n )(1 − ωq n )(1 − ω2 q n )(1 − ω3 q n )(1 − ω4 q n )
5|n 5∤n
Y Y
= (1 − q 5n )5 · (1 − q 5n )
n≥1 5∤n
Y Y Y
= (1 − q 5n )5 · (1 − q 5n ) (1 − q 5n )
n≥1 n≥1 5|n
= (q 5 ; q 5 )6∞ /(q 25 ; q 25 )∞ ,
and (2) becomes
1 (ωq; ωq)∞ (ω2 q; ω2 q)∞ (ω3 q; ω3 q)∞ (ω4 q; ω4 q)∞

= . (3)
(q; q)∞ (q 5 ; q 5 )6∞ /(q 25 ; q 25 )∞
Now we deal with the numerator of the right-hand side of (3). The pentagonal numbers
are congruent to 0, 1, or 2 modulo 5, and we can write
(q; q)∞ = 1 + q 5 − q 15 − q 35 − q 40 − q 70 · · ·
!
−q 1 − q 25 − q 50 + q 125 · · ·
!
−q 2 1 − q 5 + q 10 − q 20 − q 55 + q 75 · · · .
!
Ramanujan [4, p. 212] simply states “It can be shewn that”
(q 10 ; q 25 )∞ (q 15 ; q 25 )∞ (q 25 ; q 25 )∞
(q; q)∞ = − q(q 25 ; q 25 )∞ (4)
(q 5 ; q 25 )∞ (q 20 ; q 25 )∞
(q 5 ; q 25 )∞ (q 20 ; q 25 )∞ (q 25 ; q 25 )∞
− q2 .
(q 10 ; q 25 )∞ (q 15 ; q 25 )∞
Hardy regarded this as a gap in Ramanujan’s proof of (1), and said [2, p. 90] “Ra-
manujan never gave a complete proof.”
Proof. Here is a simple proof of (4) that I found recently. The triple product identity
can be written

∞
2 +n)/2
X
−1
(a ; q)∞ (aq; q)∞ (q; q)∞ = (−1)n a n q (n ,
n=−∞
or
2
X
(1 − a −1 )(a −1 q; q)∞ (aq; q)∞ (q; q)∞ = (−1)n (a n − a −n−1 )q (n +n)/2 .
n≥0
If we suppose a 6 = 1 and divide by 1 − a −1 , we obtain

n
a − a −n−1

(n 2 +n)/2
X
−1 n
(a q; q)∞ (aq; q)∞ (q; q)∞ = (−1) q
n≥0
1 − a −1
n+(1/2)
− a −n−(1/2)

X
n a (n 2 +n)/2
= (−1) (1/2) − a −(1/2)
q .
n≥0
a
If we now set a = exp 2iθ , we obtain

Y X sin (2n + 1)θ (n2 +n)/2
(1 − 2 cos 2θ q n + q 2n )(1 − q n ) = (−1)n q .
n≥1 n≥0
sin θ
2π
In particular, if θ = 5
,
Y
(1 + αq n + q 2n )(1 − q n ) = (q 10 ; q 25 )∞ (q 15 ; q 25 )∞ (q 25 ; q 25 )∞ (5)
n≥1
− βq(q 5 ; q 25 )∞ (q 20 ; q 25 )∞ (q 25 ; q 25 )∞ ,
while if θ = π5 ,
Y
(1 + βq n + q 2n )(1 − q n ) = (q 10 ; q 25 )∞ (q 15 ; q 25 )∞ (q 25 ; q 25 )∞ (6)
n≥1
− αq(q 5 ; q 25 )∞ (q 20 ; q 25 )∞ (q 25 ; q 25 )∞ ,
√ √
where α = 1+2 5 and β = 1−2 5 , and where we have used the triple product identity to
sum the series that arise.
By way of explanation, when we substitute θ = 2π5 , the terms of the sum in which
n ≡ 0 or −1 (mod 5) are
X
5m
sin (10m + 1) 2π5 2 +5m)/2
(−1) 2π
q (25m
m≥0
sin 5
X sin (10m − 1) 2π5 ((5m−1)2 +(5m−1))/2

+ (−1)5m−1 2π
q
m≥1
sin 5
2 2
X X
= (−1)m q (25m +5m)/2 + (−1)m q (25m −5m)/2
m≥0 m≥1
∞
2 +5m)/2
X
= (−1)m q (25m = (q 10 ; q 25 )∞ (q 15 ; q 25 )∞ (q 25 ; q 25 )∞
m=−∞

by the triple product identity. In the same way, we can sum the terms corresponding
to n ≡ 1 and n ≡ −2 (mod 5), and the terms corresponding to n ≡ 2 (mod 5) vanish.
Thus we obtain (5). We can handle (6) similarly.
If we multiply (5) and (6) together, we obtain
(q; q)∞ (q 5 ; q 5 )∞ = (q 10 ; q 25 )2∞ (q 15 ; q 25 )2∞ (q 25 ; q 25 )2∞ − q(q 5 ; q 5 )∞ (q 25 ; q 25 )∞

− q 2 (q 5 ; q 25 )2∞ (q 20 ; q 25 )2∞ (q 25 ; q 25 )2∞ ,
and if we now divide by (q 5 ; q 5 )∞ we obtain Ramanujan’s result (4).
We now complete the proof of (1). If we write
(q 10 ; q 25 )∞ (q 15 ; q 25 )∞ (q 5 ; q 25 )∞ (q 20 ; q 25 )∞
a= and b= = a −1 ,
(q 5 ; q 25 )∞ (q 20 ; q 25 )∞ 10 25
(q ; q )∞ (q ; q )∞15 25
then (4) becomes
(q; q)∞ = (q 25 ; q 25 )∞ (a − q − q 2 b)
and the numerator of the right-hand side of (3) is
(ωq; ωq)∞ (ω2 q; ω2 q)∞ (ω3 q; ω3 q)∞ (ω4 q; ω4 q)∞

= (q 25 ; q 25 )4∞
×(a − ωq − ω2 q 2 b)(a − ω2 q − ω4 q 2 b)(a − ω3 q − ω6 q 2 b)(a − ω4 q − ω8 q 2 b)
= (q 25 ; q 25 )4∞

× (a 4 − q 5 (2ab2 + b)) + q(a 3 + q 5 (ab3 + b2 )) + q 2 ((a 3 b + a 2 ) − q 5 b3 )

3 2 5 4 4 2 2
+q ((2a b + a) + q b ) + q (a b + 3ab + 1)
= (q 25 ; q 25 )4∞

4 5 3 5 2 2 2 5 3 3 5 4 4
× (a − 3q b) + q(a + 2q b ) + q (2a − q b ) + q (3a + q b ) + 5q .
So (3) becomes
X (q 25 ; q 25 )5∞
p(n)q n =
n≥0
(q 5 ; q 5 )6∞

× (a 4 − 3q 5 b) + q(a 3 + 2q 5 b2 ) + q 2 (2a 2 − q 5 b3 ) + q 3 (3a + q 5 b4 ) + 5q 4 .
If we extract those terms in which the powers are 4 modulo 5, we obtain
X (q 25 ; q 25 )5∞
p(5n + 4)q 5n+4 = 5q 4 ,
n≥0
(q 5 ; q 5 )6∞

or
X
n (q 5 ; q 5 )5∞
p(5n + 4)q = 5 ,
n≥0
(q; q)6∞
which is (1).
APPENDIX: THE TRIPLE PRODUCT IDENTITY. We have
(−aq; q 2 )∞ = (1 + aq)(−aq 3 ; q 2 )∞ .
If we write
X
(−aq; q 2 )∞ = a k ck (q)
k≥0
then c0 (q) = 1 and

X X
a k ck (q) = (1 + aq) a k q 2k ck (q).
k≥0 k≥0
It follows that for k ≥ 1,
ck (q) = q 2k ck (q) + q 2k−1 ck−1 (q),

or
q 2k−1
ck (q) = ck−1 (q).
1 − q 2k
It follows that, if we write
(a; q)k = (1 − a)(1 − aq) · · · (1 − aq k−1 ) for k ≥ 1, (a; q)0 = 1,
2
qk
ck (q) = 2 2
(q ; q )k
and
X ak q k2
2
(−aq; q )∞ = 2; q 2)
.
k≥0
(q k
This is an identity of Euler.

It is now an easy induction to prove that
∞ 2
−1 2 2
X ak q k
(−a q; q )n (−aq; q )∞ = .
k=−n
(q 2 ; q 2 )k+n
If we now let n → ∞, we obtain

∞ 2
−1 2 2
X ak q k
(−a q; q )∞ (−aq; q )∞ = 2; q 2)
,
k=−∞
(q ∞

or
∞
2
X
2 2 2 2
−1
(−a q; q )∞ (−aq; q )∞ (q ; q )∞ = ak q k ,
k=−∞
and this is the triple product identity. Note that the triple product can be written in the
equivalent forms
∞
2 +k)/2
X
3 2 3 3 3
−1
(a q; q )∞ (aq ; q )∞ (q ; q )∞ = (−1)k a k q (3k
k=−∞
and
∞
2 +k)/2
X
−1
(a ; q)∞ (aq; q)∞ (q; q)∞ = (−1)k a k q (k
k=−∞
under the substitutions (a, q) → (−aq (1/2) , q (3/2) ) and (a, q) → (−aq (1/2) , q (1/2) ) re-
spectively.
POSTSCRIPT. The “marvellous and powerful” extension of the triple product iden-
tity that I alluded to earlier is known technically as Ramanujan’s 1 ψ1 -summation.
I recommend that the reader interested in learning more about Ramanujan’s work
might begin by consulting Bruce Berndt’s book [1].
REFERENCES
1. B. C. Berndt, Number Theory in the Spirit of Ramanujan, Student Mathematical Library, vol. 34, American
Mathematical Society, Providence, RI, 2006.
2. G. H. Hardy, Ramanujan, Twelve Lectures on Subjects Suggested by His Life and Work, AMS Chelsea,
Providence, RI, 1999.
3. G. H. Hardy and E. M. Wright, An Introduction to the Theory of Numbers, Oxford University Press,
Oxford, 1960.
4. S. Ramanujan, Collected Papers, G. H. Hardy, P. V. Seshu Aiyar, B. M. Wilson, eds., AMS Chelsea,
Providence, RI, 2000.
School of Mathematics and Statistics, University of New South Wales, Sydney, Australia 2052
m.hirschhorn@unsw.edu.au

PROBLEMS AND SOLUTIONS
Edited by Gerald A. Edgar, Doug Hensley, Douglas B. West
with the collaboration of Mike Bennett, Itshak Borosh, Paul Bracken, Ezra A. Brown,
Randall Dougherty, Tamás Erdélyi, Zachary Franco, Christian Friesen, Ira M. Ges-
sel, László Lipták, Frederick W. Luttmann, Vania Mascioni, Frank B. Miles, Bog-
dan Petrenko, Richard Pfiefer, Cecil C. Rousseau, Leonard Smiley, Kenneth Stolarsky,
Richard Stong, Walter Stromquist, Daniel Ullman, Charles Vanden Eynden, Sam Van-
dervelde, and Fuzhen Zhang.
Proposed problems and solutions should be sent in duplicate to the MONTHLY

problems address on the back of the title page. Proposed problems should never
be under submission concurrently to more than one journal. Submitted solutions
should arrive before March 31, 2012. Additional information, such as general-
izations and references, is welcome. The problem number and the solver’s name
and address should appear on each solution. An asterisk (*) after the number of
a problem or a part of a problem indicates that no solution is currently available.
PROBLEMS
11600. Proposed by Michael D. Hirschhorn, University of New South Wales, Sydney,
Australia. Suppose a > 0, n ≥ 1, and 0 < r < a/n. For given θ, let

kr sin θ p
φk = arctan , ρk = a 2 − 2kra cos θ + k 2r 2 .
a − kr cos θ
Show that
Z ∞
cos(φ1 + · · · + φn ) − θ sin(φ1 + · · · + φn ) dθ π
2
= n.
0 ρ1 · · · ρn 1+θ 2a
11601. Proposed by Harm Derksen and Jeffrey Lagarias, University of Michigan, Ann
Arbor, MI. The Farey series of order n is the set of reduced rational fractions j/k in
the unit interval with denominator at most n. Let Fn be the product of these fractions,
excluding 0/1. That is,
n k−1
Y Y j
Fn = .
k=1 j=1
k
( j,k)=1
Let F n = 1/Fn . Show that F n is an integer for only finitely many n.

11602. Proposed by Roberto Tauraso, Università di Roma “Tor Vergata,” Rome, Italy.
Let p be a prime. Let Fn denote the nth Fibonacci number. Show that
X Fi
≡0 (mod p).
0<i< j<k< p
i jk

(A rational number is deemed congruent to 0 mod p if, when put in reduced form, the
numerator is a multiple of p.)
11603. Proposed by Alfonso Villani, Università di Catania, Catania, Italy. Let I be the
interval [0, ∞). Let p and r be positive, with r ≥ 1. Let f be a function on I that is
absolutely continuous in every compact interval [0, b] with b > 0. Assume that f is in
L p (I ) and that the (weak) derivative f ′ belongs to L r (I ). (Weak derivatives are part
of the theory of distributions.) Prove that limx→∞ f (x) = 0.
11604. Proposed by Pál Péter Dályay, Szeged, Hungary.
√ nGiven 0 ≤ a ≤ 2, let han i
n n
be the sequence defined by a1 = a and an+1 = 2 − 2 (2 − an ) for n ≥ 1. Find
P ∞ 2
n=1 an .
11605. Proposed by Marian Tetiva, National College “Gheorghe Roşca Codreanu,”
Bı̂rlad, Romania. Let s, R, and r be the semiperimeter, circumradius, and inradius of
a triangle with sides of length a, b, and c. Show that
√ √
R − 2r X (s − a)(s − b) X (s − c) (s − a)(s − b)
≥ −2 ,
2R c ab
and determine when equality occurs. Sums are cyclic.
11606. Proposed by Kent Holing, Trondheim, Norway. Let a, b, c, d be integers, the
first two even and the other two odd. Let Q be the polynomial x 4 + ax 3 + bx 2 + cx +
d, and assume that the Galois group of Q has order less than 24.
(a) Show that the Lagrange resolvent

3 3 2 2
x + 2b − a x
4

3 4 1
+ b − 4d + ac − a b + a x − (a 3 − 4ab + 8c)2
2 2
16 64
of Q has exactly one integer root; call it m.
(b) Show that a 2 + 4(m − b) cannot be a nonzero square.
(c) Show that if a = 0, then the Galois group of Q is cyclic if and only if
(m − b)(m + b)2 − 4c2 is square.
SOLUTIONS
An Equal Distance Sum Point

11482 [2010, 182]. Proposed by Marius Cavachi, “Ovidius” University of Constanta,
Constanta, Romania. Let n be a positive integer, and let (a1 , . . . , an ), (b1 , . . . , bn ),
and (c1 , . . . , cn ) be n-tuples of points in R2 with noncollinear centroids. For u ∈ R2 ,
let kuk be the usual euclidean norm of u. Show that there is a point p ∈ R2 such that
n
X n
X n
X
k p − ak k = k p − bk k = k p − ck k.
k=1 k=1 k=1
Solution by Robin Chapman, University of Bristol (U.K.). Define maps f a , f b , f c from

R2 to R by
n
X n
X n
X
f a ( p) = k p − ak k, f b ( p) = k p − bk k, f c ( p) = k p − ck k.
k=1 k=1 k=1
November 2011] PROBLEMS AND SOLUTIONS 847

Define F : R2 → R2 by F( p) = f a ( p) − f c ( p), f b ( p) − f c ( p) . The maps f a , f b , f c ,
and F are continuous. We must show that F has a zero (a point p such that F( p) =
(0, 0).) Suppose instead that F( p) 6= (0, 0) for all p ∈ R2 . We use the concept of
winding number. Let φ : [α, β] → R2 be a loop; that is, a continuous map with φ(α) =
φ(β). Since φ is a loop and F avoids (0, 0), F ◦ φ is a loop in R2 \ {(0, 0)}. As R2
is simply connected, φ can be deformed continuously into a path that is constant on
[α, β]; as it does, F ◦ φ deforms continuously inside R2 \ {(0, 0)} into another such
path. It follows that the winding number of F ◦ φ about (0, 0) is 0. We obtain a contra-
diction by proving that there is a loop φ such that F ◦ φ has nonzero winding number.
Define φ : [0, 2π] → R2 by
φ(θ) = (R cos θ, R sin θ ),
where R is a positive constant to be chosen later.
Let θ ∈ [0, 2π], and write p = φ(θ ). Note that k pk = R. We can take R as large
as we like; we like R > 2 maxk kak k, kbk k, kck k . Using the usual dot product in R2 ,
for 1 ≤ k ≤ n we have
1/2
kak k2

2 2 1/2
2 p · ak
k p − ak k = k pk − 2 p · ak + kak k = R 1− +
R2 R2
p · ak 2
p · ak
= R 1− 2
+ O(1/R ) = R− + O(1/R)
R R
as R → ∞, since | p · ak | = O(R). Thus
n
f a ( p) = n R − p · a + O(1/R) as R → ∞
R
where a is the centroid of the ak . Similarly
n n
f b ( p) = n R − p · b + O(1/R), f c ( p) = n R − p · c + O(1/R),
R R
where b is the centroid of the bk and c is the centroid of the ck . Hence
n
F( p) = (c − a) · p, (c − b) · p + O(1/R) as R → ∞.
R
We have assumed that a, b, c are noncollinear, so c − a and c − b are linearly inde-
pendent. Thus
c − a = (s cos α, s sin α), c − b = (t cos β, t sin β)
where s > 0, t > 0, and α − β is not an integer multiple of π . (Note that s, t, α, β are
constants and do not depend on R or θ.) Therefore

F( p) = n s cos(θ − α), t cos(θ − β) + O(1/R) as R → ∞.
Define G : [0, 2π ] → R2 by

G(θ ) := n s cos(θ − α), t cos(θ − β) .
We claim that G traces out an ellipse with center the origin exactly once, so that its
winding number about (0, 0) is ±1. To see this, let γ = α − β. Then
cos(θ − β) = cos(θ − α + γ ) = cos γ cos(θ − α) − sin γ sin(θ − α).
Now the nonsingular linear substitution (x, y) = (nsu, ntu cos γ − ntv sin γ ) con-
verts the unit circle in the (u, v)-plane into the path G. Thus G traces out an ellipse

centered at the origin, exactly once as θ increases from 0 to 2π . Hence the wind-
ing number of G is ±1. Now G(θ ) − F(φ(θ )) = O(1/R), so given ε > 0, for large
enough R we have kG(θ ) − F(φ(θ ))k < ε for all θ . Provided ε is less than the dis-
tance from the origin to the path G, the line segment from G(θ ) to F(φ(θ )) will not
meet the origin. Therefore the paths G and F ◦ φ are homotopic in R2 \ {(0, 0)}, and
so have the same winding number. This is the required contradiction.
Also solved by M. A. Prasad (India), K. Schilling, J. Simons (U. K.), R. Stong, and the proposer.
Triangle Tangent Inequality

11486 [2010, 183]. Proposed by Cezar Lupu, student, University of Bucharest,
Bucharest, Romania. Show that in an acute triangle with sides of lengths a1 , a2 , a3
and opposite angles of radian measure A1 , A2 , A3 ,
3 P3 2
( 3k=1 tan Ak )3
P
Y (1 − cos Ak ) 8 k=1 ak
≥ P3 Q3 .
k=1
cos A k 9 ( k=1 a k ) 2
k=1 (tan A k + tan A k+1 )
Solution by Pál Péter Dályay, Szeged, Hungary. Identifying A4 = A1 , in an acute tri-

angle we have 3k=1 tan Ak > 0 and 3k=1 (tan Ak + tan Ak+1 ) > 0, so the inequality
P Q
of the problem is equivalent to
Q3 P3 3
2 Y
k=1 (tan A k + tan A k+1 ) ( k=1 ak ) (1 − cos Ak ) 8
P3 3
≥ . (1)
2 cos A 9
P
( k=1 tan Ak )3 a
k=1 k k=1 k
P3 Q3
Using trigonometric formulas and the relation k=1 tan Ak = k=1 tan Ak , we obtain
Q3 ! 3 "! 3 "3 Q3
k=1 (tan A k + tan A k+1 ) Y sin(Ak + Ak+1 ) Y cos Ak cos Ak
P3 = = Q3k=1 2 .
( k=1 tan Ak )3 k=1
cos Ak cos Ak+1 k=1
sin Ak k=1 sin A k
P3 Q3 P3 (2)
2 1
Now use the relations sin A k = 4 cos(A k /2), sin A k = (3 −
P3 Q3k=1 k=1 k=1 2
k=1 cos(2A k )) = 2(1 + k=1 cos Ak ), and the sine law to obtain
( 3k=1 ak )2 ( 3k=1 sin Ak )2 16 3k=1 cos2 (Ak /2)

P P Q Q3
k=1 (1 + cos A k )
P3 2
= P 3 2
= Q 3
= Q3 . (3)
a
k=1 k k=1 sin A k 2(1 + k=1 cos A k ) 1 + k=1 cos A k
If L denotes the left-hand side of inequality (1), then using (2) and (3), we obtain
!Q " !Q " 3
3 3
cos Ak k=1 (1 + cos A k )
Y (1 − cos Ak ) 1
L = Q3k=1 2 Q3 = Q3 .
k=1 sin Ak 1 + k=1 cos Ak k=1 cos Ak 1 + k=1 cos Ak
Thus inequality (1) holds if and only if 1/(1 + 3k=1 cos Ak ) ≥ 8/9, that is, if and
Q
only if 3k=1 cos Ak ≤ 1/8. This is a known inequality, but the proof is short: note that
Q
the function f given by f (x) = log(cos x) is concave on (0, π/2) because f ′′ (x) =
− cos−2 x < 0. Thus we have
3
! ! 3 ""
X 1X 1 1
log(cos Ak ) ≤ 3 log cos Ak = 3 log = log .
k=1
3 k=1
2 8
The required inequality follows. Equality holds when the triangle is equilateral.

Also solved by A. Alt, G. Apostolopoulos (Greece), P. De (India), M. Dincǎ (Romania), J. Fabrykowski & T.
Smotzer, O. Faynshteyn (Germany), O. Kouba (Syria), O. P. Lossers (Netherlands), P. Nüesch (Switzerland),
J. Rooin (Iran), C. R. & S. Selvaraj, R. Stong, M. Tetiva (Romania), M. Vowe (Switzerland), GCHQ Problem
Solving Group (U.K.), and the proposer.
A Symmetric Inequality
11492 [2010, 278]. Proposed by Tuan Le, student, Freemont High School, Anaheim,
CA. Show that for positive a, b, and c,
√ √ √
a 3 + b3 b3 + c3 c3 + a 3 6(ab + bc + ca)
+ + ≥ √ .
a 2 + b2 b2 + c2 c2 + a 2 (a + b + c) (a + b)(b + c)(c + a)
√ √
Solution by M. A. Prasad, India. First, note that ( a 3 + b3 )/(a 2 + b2 ) ≥ 1/ a + b.
Next, put
X p
T1 = (b + c)(c + a) (a + b + c), T2 = 6(ab + bc + ca),
where the√sum is over √

all cyclic permutations
√ of {a, b, c}. It suffices to show T1 ≥ T2 .
Let x = a + b, y = b + c, z = c + a. Note z 2 < x 2 + y 2 . Also
(x 2 + y 2 + z 2 )2 x 4 + y4 + z4
ab + bc + ca = − .
4 2
Let D = 4(T1 − T2 ). Then
D = 12(x 4 + y 4 + z 4 ) + 2(x y + yz + zx)(x 2 + y 2 + z 2 ) − 6(x 2 + y 2 + z 2 )2
X X X
=3 (x 2 − y 2 )2 + 2 x y(x − y)2 − x 2 (y − z)2
X X
(x − y)2 3x 2 + 3y 2 + 8x y − z 2 ≥ (x − y)2 2x 2 + 2y 2 + 8x y ≥ 0.

=
Editorial comment. This problem can also be found at http://www.math.ust.

hk/excalibur/v14\_n2.pdf and http://ssmj.tamu.edu/problems/March-
2010.pdf with two solutions at http://www.math.ust.hk/excalibur/v14\_n3.
pdf.
Also solved by D. Beckwith, P. P. Dályay (Hungary), P. De (India), O. Faynshteyn (Germany), G. C. Greubel,
J. Grivaux (France), V. Krasniqi (Kosovo), J. H. Lindsey II, B. Mulansky (Germany), P. H. O. Pantoja (Brazil),
P. Perfetti (Italy), J. Simons (U. K.), R. Stong, L. Zhou, GCHQ Problem Solving Group (U. K.), Northwestern
University Math Problem Solving Group, and the proposer.
Glaisher–Kinkelin
11494 [2010, 279]. Proposed by Ovidiu Furdui, Campia Turzii, Cluj, Romania. Let A
be the Glaisher-Kinkelin constant, given by
n
−n 2 /2−n/2−1/12 n 2 /4
Y
A = lim n e k k = 1.2824 . . . .
n→∞
k=1
Show that
∞ (−1)n−1
Y n! A3
√ = .
n=1
2π n(n/e)n 27/12 π 1/4

Solution by Douglas B. Tyler, Raytheon, Torrance, CA. Let pn = nk=1 k k and note
Q
2 2 √
that pn ∼ An n /2+n/2+1/12 e−n /4 . Let an = n!/( 2πn(n/e)n ). By Stirling’s formula,
an tends to 1, and thus it suffices to evaluate the limit of the even-numbered partial
products. Note that
2n n n
Y k−1
Y a2k−1 Y (2k)2k
(ak )(−1) = = √ .
k=1 k=1
a2k k=1
e 2k(2k − 1)(2k − 1)2k−1
Rearranging gives
n
1 Y (2k)4k 22n(n+1) pn4
√ = √ ·
en (2n)! k=1 (2k − 1)2k−1 (2k)2k en (2n)! p2n
2 2
22n(n+1) A4 n 2n +2n+1/3 e−n A3
∼ √ · 2 2
=
(2n)n 4 4πn A(2n)2n +n+1/12 e−n π 1/4 27/12
as claimed.
Also solved by P. Bracken, B. S. Burdick, R. Chapman (U. K.), P. P. Dályay (Hungary), O. Geupel (Germany),
J. Grivaux (France), E. A. Herman, O. Kouba (Syria), V. Krasniqi (Kosovo), O. P. Lossers (Netherlands), B.
Mulansky (Germany), M. Omarjee (France), P. Perfetti (Italy), M. A. Prasad (India), P. .F. Refolio (Spain), J.
Schlosberg, J. Simons (U. K.), S. D. Smith, R. Stong, R. Tauraso (Italy), M. Tetiva (Romania), GCHQ Problem
Solving Group (U. K.), and the proposer.
Inequalities: One out of Two Ain’t Bad

11497 [2010, 370]. Proposed by Mihály Bencze, Brasov, Romania. Given n real num-
bers x1 , . . . , xn and a positive integer m, let xn+1 = x1 , and put
n n
X m X
A= xk2 − xk xk+1 + 2
xk+1 , B=3 xk2m .
k=1 k=1
Show that A ≤ 3m B and A ≤ (3m B/n)n .

Solution by M. A. Prasad, India. We first prove that
(x12 − x1 x2 + x22 )m ≤ 3m (x12m + x22m )/2. (1)
Now, −x1 x2 ≤ (x12 + x22 )/2, so (x12 − x1 x2 + x22 )m ≤ 3m (x12 + x22 )m /2m . Next,
x12(m−r ) x22r + x12r x22(m−r ) ≤ x12m + x22m , because
x12m + x22m − x12(m−r ) x22r − x12r x22(m−r ) = (x22r − x12r )(x22(m−r ) − x12(m−r ) ) ≥ 0.
Now,
m m
X m X m
2(x12 x22 )m x12(m−r ) x22r x12r x22(m−r ) x12m + x22m ,

+ = + ≤
r =0
r r =0
r

which by the binomial theorem simplifies to 2m x12m + x22m , proving (1). Therefore,
n n
X m X xk2m + xk+1
2m
A= xk2 − xk xk+1 + 2
xk+1 ≤ 3 m
= 3m−1 B.
k=1 k=1
2
The second inequality, A ≤ (3m B/n)n , is incorrect: take n = 1012 , m = 1, and x1 =

12
· · · = xn = 10−6 . Now A = 1 and B = 3, so (3m B/n)n = (9/1012 )10 < A.

Also solved by D. Beckwith, P. P. Dályay (Hungary), O. Kouba (Syria), J. H. Lindsey II, O. P. Lossers (Nether-
lands), M. A. Prasad (India), J. Simons (U. K.), R. Stong, Z. Vörös (Hungary), and GCHQ Problem Solving
Group (U. K.).
A Power Series with Nonnegative Coefficients

11501 [2010, 834]. Proposed by Finbarr Holland, University College Cork, Cork, Ire-
land. (Corrected) Let
3
g(z) = 1 − 1 1 1
.
1−az
+ 1−i z
+ 1+i z
Show that the coefficients

√ in the Taylor series expansion of g about 0 are all nonnega-
tive if and only if a ≥ 3.
Solution by Richard Stong, Center for Communications Research, San Diego, CA.
Write
∞
X az 2(a 2 − 3)z 2 + 8az 3 a 2(a 2 − 3) 2
g(z) = un zn = + = z + z + ··· .
n=1
3 3(3 − 2az + z 2 ) 3 9
√ √
Thus, if u 1 , u 2 ≥ 0, then a ≥3. If a = 3, then
√ √
3 8(z/ 3)3
g(z) = z+ √ ,
3 (1 − z/ 3)2
√
so u n = 8(n − 2)3−n/2 ≥ 0 for n ≥ 2. Now suppose a > 3 and use partial fractions
to expand in a power series as follows:
 
z 1 1 1
2
= √  √ − √ 
3 − 2az + z 2 a 2 − 3 1 − a+ a 2 −3 z a− a 2 −3
3
1− 3
z
∞
#! √ "m ! √ "m %
1 X a + a2 − 3 a − a2 − 3
= √ − zm .
2
2 a − 3 m=1 3 3
Hence all Taylor coefficients of z/(3 − 2az + z 2 ) are nonnegative. Thus

2(a 2 − 3)z 2 + 8az 3 2(a 2 − 3)z + 8az 2 z
= ·
3(3 − 2az + z 2 ) 3 3 − 2az + z 2
also has all Taylor coefficients nonnegative, and therefore g(z) does as well.
Also solved by M. Apagodu, G. Apostolopoulos (Greece), P. Bracken, N. Caro (Brazil), R. Chapman (U. K.),
D. Constales (Belgium), P. P. Dályay (Hungary), Y. Dumont (France), O. Geupel (Germany), E. A. Herman,
O. Kouba (Syria), K. McInturff, Á. Plaza & F. Perdomo (Spain), J. Simons (U. K.), M. Tetiva (Romania), M.
Vowe (Switzerland), GCHQ Problem Solving Group (U. K.), and the proposer.
A Twice Told Problem

11502, 11513 [2010, 458, 558]. Proposed by Pál Péter Dályay, Deák Ferenc High
School, Szeged, Hungary. For a triangle with area F, semiperimeter s, inradius r , cir-
cumradius R, and heights h a , h b , and h c , show that
2Fs 10r (5R − r )
5(h a + h b + h c ) ≥ + 18r ≥ .
Rr R

Solution by Vicente Vicario Garcı́a, Huelva, Spain. Let a, b, and c be the side lengths of
the triangle. We need the following result (Steinig’s inequality or the first Gerretsen in-
equality; see, for example, Bottema et. al., Geometric Inequalities, Wolters-Noordhoff,
Groningen, Netherlands, 1969, p. 50):
Lemma. In any triangle, s 2 ≥ 16Rr − 5r 2 .
This can be proved by rewriting the desired result in terms of x, y, and z, where x =
s − a, y = s − b, and z = s − c, and noting that it rearranges to be Schur’s inequality;
it can also be proved by computing that (s 2 − 16Rr + 5r 2 )/9 is the squared distance
between the incenter and centroid.
Next, rewriting Heron’s formula as
r 2 s 2 = F 2 = s(s − a)(s − b)(s − c)
= s 3 − (a + b + c)s 2 + (ab + bc + ca)s − abc
= −s 3 + (ab + bc + ca)s − abc,
and using Euler’s formula abc = 4Rr s, we obtain the (well-known) fact that in any
triangle, ab + bc + ca = s 2 + 4Rr + r 2 . From this we compute
2s 2

2Fs 1 1 1
5(h a + h b + h c ) − − 18r = 10r s + + − − 18r
Rr a b c R
10r s(ab + bc + ca) 2s 2
= − − 18r
abc R
5(s 2 + 4Rr + r 2 ) 2s 2
= − − 18r
2R R
s 2 − 16Rr + 5r 2
= ≥ 0.
2R
Also we have
2Fs 10r (5R − r ) 2s 2 10r (5R − r )
+ 18r − = + 18r −
Rr R R R
2 2
2(s − 16Rr + 5r )
= ≥ 0.
R
Thus both of the desired inequalities follow from Gerretsen’s inequality.
Editorial comment. The problem was accidentally republished as 11513.
Also solved by A. Alt, G. Apostolopoulos (Greece), M. Bataille (France), E. Braune (Austria), S. H. Brown,
B. S. Burdick, M. Can, R. Chapman (U. K.), R. Cheplyaka, V. Lucic & L. Pebody, C. Curtis, J. Fabrykowski &
T. Smotzer, O. Geupel (Germany), M. Goldenberg & M. Kaplan, B.-H. Gu (S. Korea), E. Hsynelaj (Australia)
& E. Bojaxhiu (Germany), O. Kouba (Syria), K.-W. Lau (China), J. H. Lindsey II, B. Mulansky (Germany),
P. Nüesch (Switzerland), C. R. Pranesachar (India), R. Smith, R. Stong, M. Vowe (Switzerland), Ellington
Management Problem Solving Group, GCHQ Problem Solving Group (U. K.), Mathramz Problem Solving
Group, and the proposer.

REVIEWS
Edited by Jeffrey Nunemacher
Mathematics and Computer Science, Ohio Wesleyan University, Delaware, OH 43015
Elementary Probability for Applications. By Rick Durrett. Cambridge University Press, Cam-
bridge, 2009, ix + 243 pp., ISBN 978-0-521-86756-6, $52.98.
Reviewed by Matthew Richey
Few areas of mathematics have more immediate connection to our life than proba-
bility. The past is gone, the present is ephemeral, and the future is just a guess. We
sharpen this guess through the lens of probability theory. Maybe it was easier back in
the days when we didn’t ask questions about the future, when fate was predetermined
by the gods.1 But eventually some French mathematicians got interested in the mathe-
matics of gambling. Over the next couple of centuries, probability became formalized,
axiomatized, and made profoundly useful. Unfortunately, none of this made it easier.
We have events and sample spaces; measurable (and nonmeasurable) sets; and condi-
tional, marginal, and Bayesian probabilities. We can formally define probability, but
we aren’t quite sure what it means. Is the probability associated with pulling balls out
of (Polya’s) urn the same as the probability that it will rain tomorrow? What about
those Bayesians who think of probability as belief? How can knowing the gender of
one of the twins affect the probability that the other is a girl? Even what we call “el-
ementary probability” confuses the mathematics intelligentsia. Consider, for example,
the Monte Hall problem and Marilyn vos Savant controversy [5]. No wonder students
struggle in their first probability course.
The challenge is to present probability to students in such a way that they under-
stand the fundamental concepts, appreciate the beauty and power of the subject, and—
most of all—learn to think probabilistically. There is no shortage of probability books,
ranging from texts aimed at everyone from elementary school students to professional
mathematicians, to hundreds of popular books professing to describe the uses (and
misuses) of probability theory to wider audiences (e.g., Taleb’s “Black Swan” [6]).
With all of these options available, is Durrett’s Elementary Probability for Applica-
tions really useful? Probably.
A clue to Durrett’s approach is in the title and his choice of the word “for” instead
of the usual “and.” Durrett is well acquainted with applied probability. His research
spans areas as diverse as random graphs, spatial models, and genetics, and his text
Probability: Theory and Examples [3] is standard reading for graduate students. Dur-
rett believes in the need to motivate probability through and for applications, partic-
ularly via applications that are (hopefully) relevant to the students. In this text under
review, along with the usual war horses (polling, urns, coin-flipping, and Bayes’ rule),
1 A nice description of this can be found in Mlodinow’s The Drunkard’s Walk: How Randomness Rules Our
Lives [4].

Durrett mines a number of other veins deeply: sports (especially baseball, but also
hockey and tennis), gambling (lotteries and blackjack are recurring themes), board
games, and game shows on TV (not just Monty Hall, but “Deal or No Deal” too).
There’s an interesting collection of off-beat examples—deaths from horse kicks and
shark attacks—along with topical subjects such as DNA testing.
Elementary Probability for Applications is an updated version, almost a second edi-
tion, of the author’s earlier Essentials of Probability [2], but with a greater emphasis
on minimizing the mathematical prerequisites. Elementary Probability is aimed at stu-
dents with only calculus (no linear algebra) and even those taking this course to satisfy
a general education requirement. For the most part, Durrett succeeds in bringing a rea-
sonably sophisticated approach to probability to students in the first category, less so
to students in the second category. Compared to Essentials, this text has earlier and
more intuitive introductions of key concepts such as independence, distribution, and
expected value, and copious—and often quirky—applications. Durrett’s focus is pri-
marily on discrete probability, a position made clear by his recommendation to skip
the chapter on continuous distributions and “leave this boring topic to the instructor of
the statistics course that follows yours.” There is much I like about this book. Although
I have a few quibbles here and there, they do not degrade my enthusiasm for the text.
Stylistically, the text stands out in several ways. One approach I believe the students
will appreciate is how many ideas (what we would call definitions) are presented in
context. For example, “A random variable is a numerical value determined by the
outcome of an experiment”—good enough, now show me a few. Or, “A distribution of
a discrete random variable is described by giving the value of P(X = x) for all values
of x”—excellent, now roll some dice and compute one of these distributions. Durrett
is quite generous with computations and uses numbers, graphs, and tables throughout,
always looking to put a result into context. Many authors of probability texts seem to
be satisfied with deriving a formula, perhaps assuming that readers will provide their
own calculations. For some students, certainly those for whom this text is intended,
probability needs numbers. Calculations help bring meaning to a formula. Is p bigger
than a half or less than a half? Almost zero or almost one? How does p change as n →
∞? Many examples conclude with a calculation and, even better, an interpretation.
Consider, for example, how Durrett ends the discussion of Joe DiMaggio’s 56-game
hitting streak.
To interpret our result note that the probability . . . is roughly 1/22,000 so even if
there were 220 players with 0.325 batting averages, it would take 100 years for
this to occur again.
From this, one can infer that this DiMaggio guy must have been a pretty good
ballplayer.
Another worthwhile feature of this text is Durrett’s way of interspersing more ad-
vanced mathematical concepts into elementary discussions. For example, while still
flipping coins to explain World Series winners, he mentions the chi-squared statistic,
almost as an aside. This allows Durrett to presage the notion of a p-value and hy-
pothesis testing. More sophisticated ideas, such as Stirling numbers, also pop up here
and there. There is a potential downside to this approach; topics that aren’t absolutely
necessary to the logical development of an idea could be a distraction and cause some
students to conflate key ideas with “extra for experts” concepts. Durrett trusts that the
instructor is capable of making this distinction.
The book is organized in a somewhat alternative way. Chapter One, Basic Concepts,
introduces the notion of probability via coin flipping and sports analogies. In short or-
November 2011] REVIEWS 855

der, independence, random variables, expected value, and variance (via moments) are
discussed. In this first chapter, all the positive features of the book are evident—plenty
of graphs, curiously interesting examples (e.g., an analysis of the ancient Egyptian
board game Astragali), and a steady stream of numbers. Some of the negative features
also appear. The treatment of variance is illustrative. This crucially important notion
(perhaps more interesting for statistics than probability, but still part of the core con-
tent of any probability course) is somewhat obliquely introduced within a discussion
of the expected values of functions, leading to the special case of moments, and finally
to the standard definition of variance as the expected value of the squared difference
between the random variable and the expected value of the random variable. There is
an argument for doing it this way, allowing variance to appear naturally within a more
general setting. Durrett’s presentation, in my opinion, works but potentially decreases
the probability that the student will appreciate the importance of variance.
Chapter Two, Combinatorial Probability, is where this text really hits its stride.
Durrett keeps the combinatorial complexities to a minimum; there is just enough pick-
ing and choosing to get the feel for how to use factorials and combinations. As be-
fore, everything is done in the context of examples: lotteries, card games, and the
World Series all play roles motivating the key definitions. For example, “Generalizing
from the last example we define n factorial to be. . . .” and “Passing to our general
situation, suppose we want to pick k people out of a group of n. Our first step is
to make k people standing in line. . . .” I believe that beginning mathematics students
will appreciate this approach. Potential adopters of this text will notice the absence of
the usual sequence of sections covering all the different generic discrete probability
distributions—binomial, geometric, negative bionomial, hypergeometric, etc. Durrett
develops these as needed and instead focuses on the binomial, multinomial, and Pois-
son. Instead of the menagerie of distributions, this chapter ends with sections entitled
“Card games and other urn problems,” “Probabilities of unions, Joe DiMaggio” (with
a quick mention of inclusion-exclusion and Bonferroni inequalities), and “Blackjack.”
From this, I have a theory about where Durrett picked up a lot of his applied prob-
ability. This is a fantastic chapter and place where the student can really capture the
essence of this bewitching subject.
Chapter Three, Conditional Probability, gets a little heavier. One can argue (I do)
that all of applied probability is conditional probability; the hard part is understanding
what is being conditioned upon. This (mis)understanding is one reason why students
find conditional probability so challenging. In Chapter Three the reader will con-
front some difficult probability subjects: the dreaded twin problem, Monte Hall, and
Bayesian priors and posteriors. Durrett handles all of this nicely. As usual, there are
plenty of interesting examples demonstrating and motivating the key ideas. Students
will see probability applied to the O.J. Simpson case, Simpson’s (no relation) para-
dox, mamogram posterior probabilities, and two-stage experiments. The formal topics
covered include discrete joint distributions, marginal distributions, and independence.
Near the end of the chapter, Durrett gives a wonderfully succinct explanation of the
following pleasantly surprising property of the Poisson distribution.
k
Let A1 , . . . , Ak be disjoint events whose union ∪i=1 Ai = . Suppose we per-
form the experiment a random number of times N , where N has Poisson dis-
tribution with mean λ, and let X i be the number of times Ai occurs. [Then]
X 1 , . . . , X k are independent Poissons with parameters λ P(Ai ).
This is made tangible with an example to which the students can relate.

For a concrete example, suppose that a Poisson number of cars arrive at a fast-
food restaurant each hour and let Ai be the event that the car has i passengers.
Then the number of cars with i passengers that arrive are independent Poissons.
Chapter Four, Markov Chains, is creative and artfully done, but could be challeng-
ing for students with only a calculus background. The chapter starts innocently enough
with an example-motivated introduction to Markov chains. Very quickly, however, the
level of mathematical sophistication (especially the notation) increases dramatically.
Durrett admirably works to discuss topics such as stationary distributions, limit behav-
ior of Markov chains, and absorbing chains. It will take a nimble instructor to dance
between some pretty serious theory and the practical applications of Markov chains.
For example, I question how many general education students will appreciate a theo-
rem whose short version is the sentence “the transition probability converges to equi-
librium exponentially fast.” Not having the notion of an eigenvector makes Durrett’s
discussion of stationary distributions awkward, but not unworkable. He depends heav-
ily on two- and three-state examples, supplemented by TI-83 calculator instructions.
The motivating examples are excellent: Ehrenfest chains from statistical mechanics,
the Wright-Fisher model from genetics, Monopoly (the author missed a chance to cite
the definitive work (!) on this subject [1]), inventories, and gambler’s ruin (required,
if the Blackjack section from Chapter Three is covered). Durrett deserves credit for
bringing a sophisticated approach to Markov chains to an audience with a relatively
limited mathematical background.
Next comes the aforementioned Chapter Five, Continuous Distributions. As adver-
tised, relative to what has come before, this is pretty boring and perfunctory—gone is
the spark that infuses the first four chapters. This material will be pretty intimidating
to the average calculus-only student, and completely unapproachable to students tak-
ing this course for their general education mathematics credits. No doubt, there’s a lot
of good content here: density and distribution functions, expected values, functions of
random variables, and joint/marginal/conditional distributions. The standard continu-
ous distributions, such as uniform and exponential, are discussed. The all-important
normal distribution is not mentioned in this chapter; it will appear later. Overall, the
level of discussion is extremely formal and more sophisticated than what has come
before. My advice: follow Durrett’s advice and skip this chapter.
Chapter Six, Limit Theorems, introduces the various laws of large numbers, the
normal distribution, and the central limit theorem. Arguably, all students of mathemat-
ics, whether they are taking probability as a general education requirement or on the
way to a mathematics major, deserve to encounter these fundamental ideas about ran-
domness. The level of sophistication here is, by necessity, higher than what has come
before. This chapter does not depend on knowing anything about continuous distribu-
tions and it is quite believable that a student who has been carefully and intuitively led
through the first four of Durrett’s chapters can really appreciate these advanced top-
ics. It helps that the earlier spark has returned; Durrett really wants to get these ideas
across. You can tell from the whimsical way he presents his examples: “Assuming that
dog bites are a rare event, . . . ” and “Is there a difference between baseball and flip-
ping coins?” Depending on the audience, the instructor will have to make some hard
choices in this chapter. For example, the proofs of Chebyshev’s inequality and the law
of large numbers, though short, are extremely sophisticated. I like the way the central
limit theorem is presented, but not proven, and then followed with a short section on
statistics. This chapter is worthwhile and accessible, especially if the instructor builds
up some good will beforehand.
November 2011] REVIEWS 857

The last chapter, Option Pricing, is short and sweet and also completely optional.
It might appeal to a nontrivial subset of the students in an introductory probability
course, especially those with an economics or business interest, and could serve as the
basis for student projects.
Here are a few additional observations about the text. Each chapter is augmented
with numerous exercises, organized by section, with the answers to the odd problems
provided. The index of terms is brief but there is an excellent and appreciated separate
index of examples. There are some minor typos; for example, some of the page num-
bers from the index are incorrect. In the text itself, there is a curious misspelling of
Wayne Gretzky’s last name (something this reviewer from hockey-crazed Minnesota
noticed). I have a notational quibble. Durrett uses E X to represent the expected value
of the random variable X , meaning that E X 2 represents the expected value of X 2 .
Durrett acknowledges that this can be confusing but explains, “This convention is de-
signed to cut down on parentheses.” I’m all for saving parentheses, but given their
current cost and availability, I would have appreciated a few more. I also like his con-
sistent use of “topical” examples, but do wonder if some are a little less topical now
than they were when Durrett was writing the book. For example, Wayne Gretzky and
even Roger Federer (since the rise of Rafael Nadal) are less well known now than
they were in the past. Readers of this review might recall where they were when they
heard the O. J. Simpson verdict in 1995, but users of the text probably won’t. Lastly, I
would rather an author not use the preface as a place to include a dig at the most recent
past-president of the United States.
Durrett’s Elementary Probability for Applications is an excellent addition to the
collection of probability textbooks. I’m not convinced it will work for students only
interested in satisfying a general education requirement in mathematics. But it can
work, and work well, for students with only a calculus background, especially if the
instructor is flexible and creative and approaches the subject not as a statistician or a
theoretical probabilist, but as an applied probabilist like Durrett.
REFERENCES
1. S. Abbott and M. Richey, Take a walk on the boardwalk, College Math J. 28 (1997) 162–171. http:
//dx.doi.org/10.2307/2687519
2. R. Durrett, The Essentials of Probability, Duxbury Press, Belmont, CA, 1994.
3. , Probability: Theory and Examples, 4th ed., Cambridge University Press, Cambridge, 2010.
4. L. Mlodinow, The Drunkward’s Walk: How Randomness Rules Our Lives, Pantheon, New York, 2008.
5. J. Rosenhouse, The Monty Hall Problem: The Remarkable Story of Math’s Most Contentious Brain Teaser,
Oxford University Press, New York, 2009.
6. N. Taleb, The Black Swan: The Impact of the Highly Improbable, Random House, New York, 2007.
St. Olaf College, Northfield, MN 55057

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THE AMERICAN MATHEMATICAL

The Chief Accountant and Mathematical Friend of 767

On Forming Committees 777

Complete Dissections: Converting Regions and Their 789

A “Bouquet” of Discontinuous Functions for Beginners in 799

An Introduction to Gauss Factorials 812

Enumerating the Rationals from Left to Right 830

On the Banach-Steinhaus Theorem for Topological Groups 836

Ramanujan’s “Most Beautiful Identity” 839

PROBLEMS AND SOLUTIONS 846

Elementary Probability for Applications. By Rick Durrett 854

An Official Publication of the Mathematical Association of America

Rediscovering Mathematics provides a fresh view of mathematics for

To order call 1-800-331-1622 or visit us online at www.maa.org!

http://www.editorialmanager.com/monthly Recent copies of the MONTHLY are available for pur-

2. CULTURAL BACKGROUND. We first discuss the religious and cultural back-

November 2011] S. NARAYANA AIYAR—FRIEND OF RAMANUJAN 767

3. LIFE STORY OF S. NARAYANA AIYAR. S. Narayana Aiyar was born on De-

768 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

November 2011] S. NARAYANA AIYAR—FRIEND OF RAMANUJAN 769

4. MATHEMATICAL CONTRIBUTIONS. In 1907, Narayana Aiyar became a

770 c THE MATHEMATICAL ASSOCIATION

November 2011] S. NARAYANA AIYAR—FRIEND OF RAMANUJAN 771

772 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

November 2011] S. NARAYANA AIYAR—FRIEND OF RAMANUJAN 773

Table 1. Questions by S. Narayana Aiyar.

Volume and Year No. of Proposed Questions

Progress Reports (1908) 6

774 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

November 2011] S. NARAYANA AIYAR—FRIEND OF RAMANUJAN 775

776 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

1. INTRODUCTION. As usual, the dean of your institution wishes to form a com-

November 2011] ON FORMING COMMITTEES 777

2. MAXIMUM DEGREE (LIMITED PERSONAL CONFLICT). In this section

778 c THE MATHEMATICAL ASSOCIATION

Theorem 2.2. Let G be a graph, and suppose V1 ∪ · · · ∪ Vm is a partition of V (G)

Proof. We prove Theorem 2.2 by induction on m. Let G and e = x y be as in the

November 2011] ON FORMING COMMITTEES 779

3. STRONG COLOURING (THE HAPPY DEAN PROBLEM). Suppose the

780 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

Theorem 3.1. Let G be a vertex-partitioned graph of maximum degree 1 in which

We represent a proper colouring of G with colours 1, . . . , r by a function α :

Theorem 3.2. Let G be a graph of maximum degree 1. Then sχ(G) ≤ 31.

Proof. By adding isolated vertices as necessary, we may assume n = |V (G)| is divisi-

November 2011] ON FORMING COMMITTEES 781

782 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

Theorem 4.4. (Sperner’s Lemma) Let F be a triangulation of 6m−1 with a Sperner

November 2011] ON FORMING COMMITTEES 783

Figure 5. The graph G.

784 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

November 2011] ON FORMING COMMITTEES 785

5. REMARKS. The if-and-only-if version of Theorem 4.2 is as follows.

Theorem 5.1. Let G be a graph, and suppose V1 ∪ · · · ∪ Vm is a partition of V (G)

The “only-if” implication is easy: just take I S = I ∩ V (G S ) for an IT I . The proof

786 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

ACKNOWLEDGMENTS. The author was partially supported by NSERC.

November 2011] ON FORMING COMMITTEES 787

Mathematics Is . . . and Is Not . . .

“Mathematics is abstract thought, mathematics is pure logic, mathematics is cre-

“Mathematics is not a deductive science—that’s a cliché. When you try to prove

—Submitted by Carl C. Gaither, Killeen, TX

788 c THE MATHEMATICAL ASSOCIATION OF AMERICA [Monthly 118

Dedicated to the memory of Martin Gardner