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Authors
Kaushlendra Kumar
B.Tech. (Civil Engineering), Indian Institute of Technology Kanpur
Masters in Environment Science, Jawaharlal Nehru University
Rajesh Ram
B.Tech. (Mechanical Engineering), Indian Institute of Technology Kanpur
Please e-mail: best.book4gate@gmail.com or contact @ +91 7397052994 for any suggestion or any
other information.
Engineering
Mathematics
For
GATE – 2018
CONTENT AT A GLANCE
Introduction (i)
About the Authors (ii)
Analysis (iii)
Table of Contents (iv)
Introduction
The book presents the subjects of Engineering Mathematics (as per GATE
Examination Syllabus) in a systematic, structured and precise manner from three
standpoints:
To develop their calibre, aptitude and attitude for the engineering field and
profession.
To strengthen their grasp and understanding of the concepts of the subjects
of study and their applicability at the grassroots level.
The book aims to solve problems faced by aspirants in terms of extensive
syllabus coverage, existing incongruity with syllabus, unavailability of a
standard book, etc.
Each topic in this book approaches the subject in a very conceptual and coherent
manner. While it’s illustrative and solved examples (which are taken from
previous year GATE papers from all departments) will facilitate easy
mastering of the concepts and their applications. The solved problems will
expose the students to the variety and nature of questions that they can expect to
face in the GATE examination.
This book covers all the important topics that are asked in Engineering
Mathematics for GATE examination. So this book is helpful to the GATE
aspirants who are from one of the following Engineering Streams: AE, AG, BT,
CE, CH, CS, EC, EE, IN, MA, ME, MN, MT, PE, PH, PI, TF, and XE.
Salient Features
Brief description of important theorems
Derivations of important functions, relationships and equations
Many times similar questions from previous years’ one GATE stream were
asked in another or same GATE stream. In such a case a particular question
is solved as an example and a note, like ‘Similar question was also asked
in AE-2012’ is mentioned.
From 2015 onwards, all GATE examination from all streams, numerical
answer questions apart from multiple choice questions are asked. So the
exercise questions have a mix of multiple choice questions and fill in the
blanks (in which a numerical answer is filled).
The book presents the subjects of Engineering Mathematics in a
systematic, structured and precise manner. It intends to offer GATE aspirants
a self-study and do-it-yourself approach by providing comprehensive and
step-by-step treatment of each and every aspect of the GATE examination.
In addition, the book aims to solve problems faced by aspirants in terms of
extensive syllabus coverage, existing incongruity with syllabus,
unavailability of a standard book, etc. The emphasis on fundamental
concepts helps in developing the aptitude required for success in GATE.
(i)
About the Authors
Kaushlendra Kumar is a senior Mathematics Faculty with more
than 8 years of teaching experience in various coaching institutes in
Delhi, UP and Bihar. Kumar achieved his Bachelor’s Degree in Civil
Engineering at Indian Institute of Technology Kanpur, and Master’s
degree in Environmental Sciences at Jawaharlal Nehru University.
He has made competitive material for various GATE and also for
JEE (Mains & Advance) Coaching Institutes. Many of the GATE
aspirants, after getting guidance by Kumar, were able to secure
SINGLE and TWO DIGIT Ranks in GATE Examination. In the
past two years, he has focused on developing this book in such a
ways that he is going to appear in GATE 2017 examination.
(ii)
Stream wise analysis of weightage (out of 100 marks) of
Engineering Mathematics in Previous Year GATE Examinations
Year/ AE AG BT CE CH CS
Sets Set 1 Set 1 Set 1 Set 1 Set 2 Set 1 Set 1 Set 2 Set 3
2010 14 15 5 15 - 16 16 - -
2011 11 13 6 13 - 13 14 - -
2012 13 13 6 13 - 13 16 - -
2013 12 9 13 8 - 12 13 - -
2014 15 13 13 13 13 14 13 13 13
2015 13 13 13 13 13 13 11 11 13
2016 13 13 10 14 15 13 11 11 -
2017 15 13 15 13 12 13 10 13 -
Average 13.25 12.75 10.125 12.92 13.375 12.71
Year/ EE EC IN XE
Sets Set 1 Set 2 Set 3 Set 1 Set 2 Set 3 Set 4 Set 1 Set 1
2010 13 - - 13 - - - 12 15
2011 11 - - 8 - - - 10 15
2012 10 - - 11 - - - 15 15
2013 13 - - 10 - - - 10 15
2014 17 11 14 13 13 13 13 10 15
2015 12 12 - 15 14 15 - 10 15
2016 16 16 - 13 13 13 - 15 15
2017 16 13 - 15 15 - - 13 15
Average 13.38 12.94 11.875 15
Year/ ME MN MT PE PH PI TF
Sets Set 1 Set 2 Set 3 Set 4 Set 1 Set 1 Set 1 Set 1 Set 1 Set 1
2010 9 - - - 16 15 - 9 17 15
2011 13 - - - 13 10 - 9 15 9
2012 13 - - - 9 13 - 8 15 10
2013 14 - - - 8 14 - 10 15 8
2014 13 13 13 13 14 11 - 7 10 14
2015 15 15 15 - 15 13 - 6 13 13
2016 13 13 13 - 12 13 15 9 13 14
2017 13 13 - - 12 13 25 10 13 13
Average 13.19 12.375 12.75 20 8.5 13.875 12
(iii)
Table of Content
Engineering Mathematics
Page Number
List of Chapters and their sub-topics
From - To
GATE - 2017: Chapter - 0: Prerequisite Chapter - 0 [1] - Chapter - 0 [2]
Chapter 0: Prerequisite [0.1] - [0.60]
0.1 Equations and In-equations [0.1] - [0.15]
0.2 Sequence, Series and Progression [0.16] - [0.20]
0.3 Binomial Theorem [0.20] - [0.22]
0.4 Trigonometry [0.23] - [0.33]
0.4 Two Dimensional Geometry [0.33] - [0.45]
0.5 Three Dimensional Geometry [0.46] - [0.59]
GATE - 2017: Chapter - 1: Linear Algebra Chapter - 1 [1] - Chapter - 1 [9]
Chapter 1: Linear Algebra [1.1] - [1.103]
1.1 Matrices and Determinants [1.1] - [1.27]
1.1.1 Algebra of Matrices [1.2] - [1.5]
1.1.2 Properties of Matrices [1.6] - [1.8]
1.1.3 Special Matrices [1.8] - [1.10]
1.1.4 Determinant of a Square Matrix [1.10] - [1.13]
1.1.5 Operations in Determinants [1.14] - [1.14]
1.1.6 Properties of Determinants [1.14] - [1.18]
1.1.7 Adjoint and Inverse of a Square Matrix [1.18] - [1.24]
1.2 System of Linear Equations [1.28] - [1.40]
Solution of a Non-Homogeneous System of Linear
1.2.1 [1.28] - [1.37]
Equations
Solution of Homogeneous Linear Equations
1.2.2 [1.38] - [1.39]
Ax = O
1.3 Vector Algebra [1.41] - [1.52]
1.3.1 Scalar or Inner or Dot Product of two Vectors [1.44] - [1.45]
1.3.2 Vector or Cross Product of two Vectors [1.46] - [1.48]
1.3.3 Scalar Triple Product [1.48] - [1.49]
1.3.4 Vector Triple Product [1.50] - [1.50]
Vector Space, Linear Dependence and Rank of a
1.4 [1.53] - [1.74]
Matrix
1.4.1 Vector Space and Sub-Space [1.53] - [1.54]
1.4.2 Linear Dependence, Spanning Sets and Bases [1.54] - [1.57]
1.4.3 Dot – Product, Norm and Angle for n – Space [1.57] - [1.63]
1.4.4 Rank of a Matrix [1.63] - [1.72]
Eigenvalue (Characteristic root) and Eigenvector
1.5 [1.75] - [1.103]
(Characteristic Vector)
1.5.1 Eigenbases of a Matrix [1.91] - [1.92]
1.5.2 Similar Matrices [1.92] - [1.96]
1.5.3 Matrix Polynomial [1.96] - [1.99]
1.5.4 Quadratic Forms [1.99] - [1.100]
GATE - 2017: Chapter - 2: Calculus Chapter - 2 [1] - Chapter - 2 [11]
Chapter 2: Calculus [2.1] - [2.167]
2.1 Set Theory [2.1] - [2.11]
2.1.1 Venn – Euler Diagram [2.2] - [2.4]
2.1.2 Relations [2.5] - [2.8]
(iv)
2.2 Function [2.11] - [2.21]
2.2.1Types of Functions [2.11] - [2.15]
2.2.2Classification of Function [2.16] - [2.19]
2.2.3Composite Function [2.19] - [2.20]
Limit, Continuity and Differentiability of a
2.3 [2.22] - [2.48]
Function
2.3.1 Limit of a Function [2.22] - [2.30]
2.3.2 Continuity of a Function [2.31] - [2.35]
2.3.3 Differentiability of a Function [2.35] - [2.45]
2.4 Application of Derivatives [2.49] - [2.75]
2.4.1 Tangent and Normal [2.49] - [2.51]
2.4.2 Monotonicity of a Function [2.51] - [2.64]
2.4.3 Mean Value Theorem [2.64] - [2.67]
2.4.4 Partial and Total Derivatives of a Function [2.67] - [2.74]
2.5 Integration [2.76] - [2.105]
2.5.1 Indefinite Integrals [2.76] - [2.83]
2.5.2 Definite Integrals [2.83] - [2.96]
Curve Sketching, Arc Length, Area and Volume of
2.5.3 [2.96] - [2.103]
a Curve
2.6 Vector Differential Calculus [2.106] - [2.120]
2.6.1 Scalar and Vector Fields [2.106] - [2.108]
2.6.2 Derivative of a Scalar Field [2.108] - [2.113]
2.6.3 Derivative of a Vector field [2.113] - [2.118]
2.6.4 Vector/Matrix Derivatives [2.118] - [2.119]
2.7 Vector Integral Calculus [2.121] - [2.147]
2.7.1 Line Integrals [2.121] - [2.126]
2.7.2 Double Integrals [2.127] - [2.131]
2.7.3 Green’s Theorem in a Plane [2.131] - [2.133]
2.7.4 Surface Integrals [2.134] - [2.138]
2.7.5 Volume Integrals [2.138] - [2.140]
2.7.6 Gauss’s Divergence Theorem [2.140] - [2.143]
2.7.7 Stokes’ Theorem [2.143] - [2.144]
2.8 Sequence and Series [2.148] - [2.165]
2.8.1 Sequence [2.148] - [2.149]
2.8.2 Series [2.149] - [2.154]
2.8.3 Power Series [2.154] - [2.156]
2.8.4 Taylor’s Series and Maclaurin Series [2.156] - [2.159]
2.8.5 Improper Integral [2.159] - [2.163]
GATE - 2017: Chapter - 3: Differential Equation Chapter - 3 [1] - Chapter - 3 [9]
Chapter 3: Differential Equation [3.1] - [3.104]
3.1 Solution of a Differential Equation [3.4] - [3.19]
Solution of first order and first degree differential
3.1.1 [3.4] - [3.13]
equation
3.1.2 Exact Differential equation [3.13] - [3.16]
Solution of Special type of second order differential
3.1.3 [3.16] - [3.18]
equation
3.2 Higher Order Differential Equation [3.20] - [3.44]
3.2.1 Preliminary theory: Linear Equations [3.19] - [3.22]
3.2.2 Reduction of Order [3.22] - [3.23]
(v)
Homogeneous Linear Equations with Constant
3.2.3 [3.24] - [3.30]
Coefficients
Non – Homogeneous Linear Equation with
3.2.4 [3.30] - [3.36]
Constant coefficients
3.2.5 Variation of Parameters [3.36] - [3.39]
3.2.6 Cauchy – Euler Equation [3.39] - [3.44]
3.3 Laplace Transform [3.47] - [3.70]
3.3.1 Inverse Transforms and Transforms of Derivatives [3.49] - [3.54]
3.3.2 Important Singularity Functions [3.54] - [3.57]
3.3.3 Operational Properties [3.57] - [3.68]
3.4 Fourier Series [3.71] - [3.86]
3.4.1 Waveform Symmetry [3.72] - [3.76]
3.4.2 Convergence and Sum of a Fourier Series [3.76] - [3.78]
3.4.3 Half Range Expansion [3.79] - [3.79]
3.4.4 Complex Form of Fourier Series [3.79] - [3.80]
3.4.5 Fourier Transform [3.80] - [3.84]
3.5 Partial Differential Equation [3.87] - [3.104]
Modelling and Solution of Vibrating String: Wave
3.5.1 [3.93] - [3.96]
Equation
3.5.2 Modelling and Solution of Heat Equation [3.96] - [3.99]
Solution of Steady Two – Dimensional Heat
3.5.3 [3.99] - [3.102]
Problem: Laplace’s Equation
GATE - 2017: Chapter - 4: Complex Variables Chapter - 4 [1] - Chapter - 4 [5]
Chapter - 4: Complex Variables [4.1] - [4.38]
4.1 Basic Concept of Complex Numbers [4.1] - [4.17]
4.1.1 Complex Plane [4.2] - [4.3]
4.1.2 Polar Form of a Complex Number [4.4] - [4.7]
4.1.3 Complex Function [4.7] - [4.15]
4.1.4 Logarithm of a Complex number [4.15] - [4.16]
4.2 Complex Integration [4.17] - [4.35]
4.2.1 Line Integral in the Complex Plane [4.17] - [4.18]
4.2.2 Cauchy’s Integral Theorems [4.18] - [4.21]
4.2.3 Derivatives of analytic function [4.21] - [4.21]
4.2.4 Sequence and Series [4.22] - [4.25]
4.2.5 Singularities, Zero and Infinity [4.25] - [4.26]
4.2.6 Zeros of Analytic function [4.26] - [4.26]
4.2.7 Residue Integration Method [4.26] - [4.35]
GATE - 2017: Chapter - 5: Numerical Methods Chapter - 5 [1] - Chapter - 5 [5]
Chapter - 5: Numerical Methods [5.1] - [5.56]
5.1 Algebraic and Transcendental Equation [5.3] - [5.26]
5.1.1 Location of real roots of an equation [5.3] - [5.3]
Exact Solution of Algebraic and Transcendental
5.1.2 [5.3] - [5.4]
Equation
Numerical Solution of Algebraic and
5.1.3 [5.4] - [5.16]
Transcendental Equation
5.1.4 Rate of Convergence of Approximate Methods [5.16] - [5.16]
5.1.5 Interpolation [5.16] - [5.23]
(vi)
Numerical Solution of Definite Integrals, First
5.2 [5.27] - [5.56]
Order ODEs, and System of Linear Equations
GATE - 2017: Chapter - 6: Probability and Statistics Chapter - 6 [1] - Chapter - 6 [6]
(vii)
Engineering Mathematics Chapter – 0: Prerequisite Chapter – 0 [1]
[BT – 2017 (2 marks)]: Consider an infinite number of cylinders. The first cylinder has a radius of 1
meter and height of 1 meter. The second one has a radius of 0.5 meter and height of 0.5 meter. Every
subsequent cylinder has half the radius and half the height of the preceding cylinder. The sum of the
volumes (in cubic meters) of these infinite number of cylinders is _____. Given data: 3.14 .
Solution: The required sum of the volumes (in cubic meters) is
Solution: Let r and h be the radius and height, respectively, of the cylinder. So volume of the
cylinder is V r 2 h . As per data given r h ; so V r 3 . Thus V1 (1)3 , V2 (1 2)3 ;
V3 (1 4)3 , …. Hence total volume V V1 V2 V3 {13 (1 2)3 (1 4)3 } . The infinite
series 13 (1 2)3 (1 4)3 is in GP whose first term is 1 and common ratio is (1 2)3 , hence the
sum of the infinite series is 1 {1 (1 2)3 } 8 7 . So total volume is V (8 7) 3.59 m3.
[EC – 2017 (1 mark)]: The smaller angle (in degrees) between the planes x y z 1 and
2 x y 2 z 0 is _____.
Solution: The normal vector to the given planes is given by: n1 i j k and n 2 2i j 2k . So
smaller angle ( ), i.e. acute angle between the given planes the same as the acute angle between the
n1 n 2 1 2 1 2
normal vectors. So cos 1 cos
1
cos (0.5773) 54.73
o
n1 n 2 3 9
[EE – 2017 (1 mark)]: Let x and y be integers satisfying the following equations 2 x 2 y 2 34 ;
x 2 y 11 . The value of ( x y ) is _____.
Solution: Substituting 2nd equation in 1st equation, we get 2 x 2 y 2 34 2(11 2 y ) 2 y 2 34
9 y 2 88 y 208 0 y 5.77 or y 4 . As y is an integer so we take y 4 . On
substituting y 4 in 2nd equation we get, x 3 ; thus x y 3 4 7 .
[MN – 2017 (1
mark)]: Which one of
the following plots
represents the
relationship xy c ,
where c is positive.
(a) I (b) II (c) III (d) IV
Solution (c): As xy c , where c is positive; so xy 0 , which is only possible is either both x , y
are positive or both x , y are negative. This condition is represented in figure III. So option (c) is
correct.
Chapter 0 : Prerequisite
0.1 Equations and In-equations
A set of all real numbers can be expressed as: x ( , ) or x R or x
x ( , a) (b, ) x R [a, b] x ( , a] [b, ) x R (a, b)
x (, a ) [b, ) x R [ a, b) x ( , a] (b, ) x R ( a, b]
Intervals: The set of numbers between any two real numbers is called
interval. There are three four types of interval:
Closed interval: x [a, b] or { x : a x b} , as shown in Fig. 0.1(a).
Open interval: x ( a, b) or { x : a x b} , as shown in Fig. 0.1(b).
Open-Closed interval: x ( a, b] or { x : a x b} , as shown in Fig.
0.1(c).
Closed-Open interval: x [ a, b) or { x : a x b} , as shown in Fig.
0.1(d).
Figure 0.1: Types of Intervals
Inequalities
If a and b are real and (i) ( a b) is positive then a b (ii) ( a b) is negative then a b
If c is real and a b then a c b c and a c b c
If m is real and a b , then
am bm if m 0
am bm if m 0
m m m m
a b if m 0
a b if m 0
a a ax a a ax
For any x 0 If 0 1 , then 1 If 1 , then 1
b b bx b b bx
x (1 x ) 2 for all x 0 , equality holds at x 1
x (1 x) 2 for all x 0 , equality holds at x 1 .
( a12 a22 an2 )(b12 b22 bn2 ) ( a1b1 a2 b2 an bn ) 2
Squaring an inequality: If a b , then it is not necessary that a 2 b 2 always holds true. For e.g.
2 2 2 2 2 2
2 5 2 5 ; but 5 2 (5) 2 . Thus a b a b follows only when a b .
Law of reciprocal: If both sides of inequality have same sign, then after taking their reciprocal
the sign of inequality gets reversed, i.e. if 0 a b (1 b ) (1 a ) ; or if
a b 0 (1 b ) (1 a ) . If both sides of inequality have opposite sign, then after taking their
reciprocal the sign of inequality remains same, i.e. if a 0 b (1 a ) (1 b ) .
2. The solution of f ( x ) 0 is the union of all intervals in which we have put the plus sign and the
solution of f ( x ) 0 is the union of all intervals in which we have put the minus sign.
It is to be noted that, we put plus sign in the interval of the right of the largest of these numbers, i.e. on
the right of an , only when all the ‘ x ’ in (i) are positive. If some of the ‘ x ’ in (i) are negative then
Put plus sign in the interval of the right of largest of these numbers, i.e. on the right of an , if total
number of ve ‘ x ’ are even. Then follow the rules given in 1 and 2 for finding the answer.
Put minus sign in the interval of the right of largest of these numbers, i.e. on the right of an , if
total number of ve ‘ x ’ are odd. If kn is even then put minus sign on left of an and if kn is odd
then put plus sign on the left of an . Then follow the rules given in 1 and 2 for finding the answer.
Example 0.3 [TF-2014 (1 mark)]: The range of values of x satisfying the inequality x 2 3 x 2 0
is
(a) x 0 (b) 0 x 1 (c) 1 x 2 (d) cannot be determined
2
Solution (c): x 3 x 2 0 ( x 2)( x 1) 0 , now equating each factor to ‘0’
and from wavy curve method, we have, x (1, 2) or 1 x 2 .
Polynomial Equation: Algebraic expression containing terms of the form cx n , where n being a
non-negative integer, is called a polynomial. A polynomial of degree ‘ n ’ can be represented as:
f ( x ) an x n a n 1 x n 1 an 1 x n 1 a2 x 2 a1 x a0 , where x is a variable, a0 , a1 , a2 , , an are
constants and an 0 . A polynomial is called to be real if the constants and variable are real; and is
called to be complex if constants or variable or both are complex numbers. A polynomial of second
degree is generally called a quadratic polynomial. Polynomials of degree 3 and 4 are known as cubic
and biquadratic polynomials respectively. If f ( x ) is a polynomial, real or complex, then f ( x ) 0 is
called a polynomial equation. The values of variable ‘ x ’ satisfying a given polynomial equation are
called its roots, i.e. if x is a root of equation f ( x ) 0 , then f ( ) 0 .
Descarte’s rule of sign for finding number of roots of a polynomial equation: Let f ( x ) 0 be a
polynomial equation with real coefficients,
The number of ve roots of f ( x ) 0 is equal to the number of variation of signs of the
successive coefficients of f ( x ) or less by an even number.
The number of ve roots of f ( x ) 0 is equal to the number of variation of signs of the
successive coefficients of f ( x ) or less by an even number.
Let f ( x ) 0 have no ‘zero’ root, and let the number of variations of signs of the successive
coefficients of f ( x ) and f ( x ) be r and k , respectively, then
The maximum number of real roots of f ( x ) is ( r k )
If ( r k ) n , then the equation f ( x ) 0 has exactly ‘ r ’ number of ve real roots and ‘ k ’
number of ve real roots.
Example 0.4: Find the possible number of roots of the polynomial equation f ( x ) 0 , where
f ( x ) 3 x5 2 x 4 5 x 3 4 x 2 6 x 7 .
Solution: f ( x ) 3 x 5 2 x 4 5 x 3 4 x 2 6 x 7 two change of ve ve
complex total
signs, so f ( x ) 0 has either two ve real roots or ‘zero’ ve real real real
2 3 0 5
roots. f ( x ) 3x 5 2 x 4 5 x 3 4 x 2 6 x 7 three change of 2 1 2 5
sign, so f ( x ) 0 has either three ve real roots or one ve real 0 3 2 5
root and two complex roots. The possibility of roots is given in 0 1 4 5
table, in which any one row is possible.
Example 0.5: Find how many ve real roots lie in (0,1) and (1, 2) for the polynomial equation
f ( x ) 27 x 4 48 x 2 12 x 13 0 .
Solution: f ( x ) 27 x 4 48 x 2 12 x 13 0 two variation of signs, so f ( x ) 0 has two ve real
roots or no ve real roots. Now f (0) 13 0 , f (1) 20 0 and f (2) 229 0 ; thus f ( x ) 0
has at least one ve real root in (0,1) and at least one ve real root in (1, 2) . Thus, we can say that
f ( x ) 0 has exactly one ve real root in (0,1) and exactly one ve real root in (1, 2) .
Example 0.6 [EE-1994 (2 marks)]: Find the number of positive real roots of the equation
3
x 2x 2 0 .
Solution: As f ( x) x 3 2 x 2 we have two variation of sign so f ( x ) 0 has two ve real
roots or no ve real root. Also f ( x ) x 3 2 x 2 we have one variation
of sign so f ( x) 0 has one ve real root. Now
f ( x ) 3 x 2 2 0 x 2 3 ; also as x , f ( x ) and as
x , f ( x) . Since f (0) 2 and f
2 3 0 , so from x 0 to
x 2 3 , the graph of f ( x ) 0 not crosses the x axis. Hence the possible graph is shown in
figure, which shows that f ( x ) 0 has no ve real roots and thus f ( x ) 0 has two complex roots
and one ve real root. Thus f ( x ) 0 has no ve real roots.
Example 0.7 [CS-1997 (2 marks)]: A polynomial p ( x ) is such that p (0) 5 (a) 1 (b) 2
, p (1) 4 , p (2) 9 and p (3) 20 . The minimum degree it can have is (c) 3 (d) 4
Solution (b): As x approaching from x 0 to x 1 , the value of p( x ) decreases; on the other hand
as x approaching from x 1 to x 2 to x 3 , the value of p( x ) keeps on increasing; so from the
given conditions, we can say that between x 1 and x 2 , the graph of p( x ) either cut the x axis
two times (i.e. it has two real roots) or it takes the shape of concave downward with cutting the x
axis (i.e. it has two complex roots); in either case we have minimum two roots.
Example 0.10 [CE-2007 (2 marks)]: Given that one root of the equation x 3 10 x 2 31x 30 0 is
5, the other two roots are
Example 0.12 [IN-2008 (2 marks)]: It is known that two roots of the non-linear equation
3 2
x 6 x 11x 6 0 are 1 and 3. The third root will be
(a) j (b) j (c) 2 (d) 4
Solution (c): Let the third root be . As the sum of roots be (6 1) 6 1 3 6 2 .
Example 0.16 [XE-2014 (1 mark)]: If a cubic polynomial passes through the points (0,1) , (1, 0) ,
(2,1) and (3,10) , then it also passes through the point
(a) ( 2, 11) (b) (1, 2) (c) (1, 4) (d) ( 2, 23)
Solution (b): Let the cubic polynomial is f ( x ) ax 3 bx 2 cx d . Now satisfying the given point,
we get: (0,1) d 1 f ( x ) ax 3 bx 2 cx 1 …(i); (1, 0) a b c 1 …(ii);
(2,1) 8a 4b 2c 0 …(iii); (3,10) 27 a 9b 3c 9 …(iv). From (ii), (iii) and (iv), we get
3 2
a 1 , b 2 and c 0 . So f ( x ) x 2 x 1 , now checking all the option we will find that the
point ( 1, 2) satisfy f ( x ) x 3 2 x 2 1 , so option (b) is correct.
(a) (b)
(c)
(d)
(e) (f)
Figure 0.3: Condition for position of roots of Quadratic Equation (Eq. 0.1)
Condition for a number k , if both the roots of Eq. 0.1 are less than k : (i) D 0 (roots
may be equal) (ii) a f ( k ) 0 (iii) k b 2a , where ; as shown in Fig. 0.3(a).
Condition for a number k , if both the roots of Eq. 0.1 are greater than k : (i) D 0
(roots may be equal) (ii) a f ( k ) 0 (iii) k b 2a , where ; as shown in Fig.
0.3(b).
Condition for a number k , if k lies between the roots of Eq. 0.1: (i) D 0 (ii)
a f (k ) 0 , where ; as shown in Fig. 0.3(c).
Condition for numbers k1 and k 2 , if exactly one root of Eq. 0.1 lies in the interval ( k1 ,
k 2 ): (i) D 0 (ii) f ( k1 ) f (k 2 ) 0 , where ; as shown in Fig. 0.3(d).
Condition for numbers k1 and k 2 , if both roots of Eq. 0.1 are confined between k1 and
k 2 : (i) D 0 (roots may be equal) (ii) a f ( k1 ) 0 (iii) a f ( k2 ) 0 (iv)
k1 b 2a k 2 , where and k1 k2 ; as shown in Fig. 0.3(e).
Condition for numbers k1 and k 2 , if k1 and k 2 lie between the roots of Eq. 0.1: (i)
D 0 (ii) a f ( k1 ) 0 (iii) a f ( k2 ) 0 , where ; as shown in Fig. 0.3(f).
x, if x 0
Modulus or Absolute value function: The function defined by f ( x) x 0, if x 0 is called
x, if x 0
the modulus function. Its Domain: ; Range: [0, ) ; Period: Non – periodic; Nature: Even; Interval
in which the inverse can be obtained: ( , 0] [0, ) .
Please note that the questions on Equations and In-equations also came in previous years
‘General Aptitude’ section in GATE examination. These questions are given in ‘General
Aptitude’ section.
Exercise 0.1
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
21.
If x is real, then greatest and least values of ( x 2 x 1) ( x 2 x 1) are
(a) 3, 1 2 (b) 3,1 3 (c) 3, 1 3 (d) None of these
22. If f ( x ) is quadratic expression which is positive for all real value of x and
g ( x) f ( x) f ( x ) f ( x ) . Then for any real value of x
(a) g ( x ) 0 (b) g ( x ) 0 (c) g ( x ) 0 (d) g ( x ) 0
23. If , ( ) are roots of the equation x 2 bx c 0 where (c 0 b) then
(a) 0 (b) 0 (c) 0 (d) 0
24. If the roots of the equation x 2 2ax a 2 a 3 0 are real and less than 3, then
(a) a 2 (b) 2 a 3 (c) 3 a 4 (d) a 4
25. The value of a for which 2 x 2 2(2a 1) x a ( a 1) 0 may have one root less than a and
another root greater than a are given by
(a) 1 a 0 (b) 1 a 0 (c) a 0 (d) a 0 or a 1
26. 5 2
The maximum possible number of real roots of equation x 6 x 4 x 5 0 is _____.
27. In solving a problem, a person makes a mistake in the coefficient of the first degree term and
obtains –9 and –1 for the roots. Another person makes a mistake in the constant term of the
equation and obtains 8 and 2 for the roots. The correct quadratic equation was
(a) x 2 10 x 9 0 (b) x 2 10 x 16 0 (c) x 2 10 x 9 0 (d) None of these
28. The value of x 2 4 x 11 0 can never be less than _____.
29. If the roots of the equation ( a 2 bc ) x 2 2(b 2 ac ) x (c 2 ab) 0 are equal, where b 0 ,
then which one of the following is correct?
(a) a b c abc (b) a 2 b 2 c 2 0
(c) a3 b3 c3 0 (d) a3 b3 c 3 3abc
30. If the roots of the equation ax 2 bx c 0 are 1 and 1, then which one of the following is
correct?
(a) a and c are both zero (b) a and b are both positive
(c) a and c are both negative (d) a and c are of opposite sign
31.
If x 2 6 6 6 6 , then what is positive value of x ? _____
32. The difference in the roots of the equation 2 x 2 11x 5 0 is _____.
33. The quadratic equation whose roots are 3 and –1 is
(a) x 2 4 x 3 0 (b) x 2 2 x 3 0 (c) x 2 2 x 3 0 (d) x 2 4 x 3 0
34. If one root of the equation (1 a) x 2 (1 b ) x (1 c ) 0 is reciprocal of the other, then which
one of the following is correct?
(a) a b (b) b c (c) ac 1 (d) a c
35. 2 2
If the roots of the equation x 2ax a a 3 0 are real and less than 3, then which one of
the following is correct?
(a) a 2 (b) 2 a 3 (c) 3 a 4 (d) a 4
36. 2
Two students A and B solve an equation of the form x px q 0 . A starts with a wrong
value of p and obtains the roots as 2 and 6. B starts with a wrong value of q and gets the
roots as 2 and –9. What is the product of roots of the correct equation? _____
37. If one of the roots of 7 x 2 50 x k 0 is 7, then the value of k is _____.
38. If sin and cos are the roots of the equation ax 2 bx c 0 , then which one of the
following is correct?
(a) a 2 b 2 2ac 0 (b) a 2 b 2 2ac 0
(c) a 2 c 2 2ab 0 (d) a 2 b 2 2ac 0
39. What is the condition that the equation ax 2 bx c 0 , where a 0 has both the roots
positive?
(a) a , b and c are of same sign (b) b and c have the same sign opposite to that of a
(c) a and b are of same sign (d) a and c have the same sign opposite to that of b
40. The equation (1 n ) x 2ncx (c a 2 ) 0 will have equal roots if
2 2 2
_____
47. What is the least integral value of k for which the equation x 2 2( k 1) x (2k 1) 0 has
both the roots positive? _____
48. When the roots of the quadratic equation ax 2 bx c 0 are negative of reciprocals of each
other, then which one of the following is correct?
(a) b 0 (b) c 0 (c) a c (d) a c
49. 2
For what value of k , will the roots of the equation kx 5 x 6 0 be in the ratio of 2 : 3 ?
_____
50. Students of a class are made to stand in rows. If one student is extra in a row, there would be
two rows less. If one students is less in a row, there would be three rows more. Then, what is
the number of students in the class? _____
51. x( x 1) ( m 1) x
If the roots of the equation are equal, then what is the value of m ?
( x 1)( m 1) m
(a) 1 (b) 0 (c) 1 2 (d) 1 2
52. What is one of the roots of the equation 2 x (3 x) (3 x ) 2 x 3 2 ?
(a) 1 (b) 2 (c) 3 (d) 4
53. 2
If and are the roots of the equation x 3 x 2 0 , then which equation has the roots
( 1) and ( 1) ?
(a) x 2 5 x 6 0 (b) x 2 5 x 6 0 (c) x 2 5 x 6 0 (d) x 2 5 x 6 0
54. For what value of k , does the equation kx 2 (2k 6) x 16 0 have equal roots?
(a) 1 and 9 (b) –9 and 1 (c) –1 and 9 (d) –1 and –9
55. 2 2
If the product of the roots of x 3kx 2k 1 0 is 7 for a fixed k , then what is the nature
of the roots?
(a) integral and positive (b) integral and negative
(c) irrational (d) rational but not integral
56. 2 2
If x kx 21 0 and x 3kx 35 0 have one common root, then sum of all the values of
k is _____.
(a) x ( , 2] {1} (b) x [2, ) {1} (c) x [2, ) {1} (d) x (, 2] {1}
87. 2 5 4
If x( x 2) ( x 1) (2 x 3)( x 3) 0 , then
(a) x [0,1 2] [3, ) (b) x [0,1) [3 2 , )
(c) x [0,1 2] [3, ) (d) x [0,1] [3 2, )
88.
Find the range of values of y if y ( x 2 1) ( x 2 2)
(a) y 1 2 (b) y 1 (c) y 1 2 or y 1 (d) y 1 2 or y 1
89. If x 1 3 x , then
(a) 2 x 3 (b) 2 x 3 (c) 2 x 3 (d) 2 x 3
90. If x 2 x 2 x 3 , then
(a) x 0 (b) x 0 (c) x 3 (d) x 3
91. 2 2
If ( x 4) x 1 0 , then
(a) x (2,1) (1, 2) (b) x (2, 1] [1, 2) (c) x ( 2, 1] (d) x [1, 2)
92. 2
If 2 x 1 x 1 3 , then the sum of values of x is
(a) 1 (b) 2 (c) –2 (d) –1
If a constant is added to (or subtracted from) every term of an AP, the resulting sequence will also
be an AP with same common difference.
If every term of an AP (with common difference d ) is multiplied by (or divided by) a fixed
constant ( k ), the resulting sequence will also be an AP with common difference dk (in case of
multiplication) and d k in case of division.
If respective terms of two AP are added or subtracted, the new progression also in AP.
The sum of the term which are equidistance from the beginning to the end is constant, and is
equal to the sum of first and last term.
Example 0.19 [TF-2014 (2 marks)]: The 51st common term of the two arithmetic sequence 15, 19,
23, … and 14, 19, 24, …, is _____
Solution: The 1st AP: 15, 19, 23, 27, 31, 35, 39, 43, ; The 2nd AP is: 14, 19, 24, 29, 34, 39, . So
the arithmetic sequence of common terms between the given two APs is 19, 39, . Thus 51st
common term is t51 19 (51 1)20 1019 .
Geometric Progression: If the terms of a series increases or decreases by a common ratio then the
2
series is called a geometric progression (GP). The general form of a GP is a, ar , ar , , where a is
the first term and r is the common ratio. The nth term of a GP is, t n ar n 1
a ( r n 1)
Sum of n terms of the GP is, S n for all r 1 .
( r 1)
Sum of an infinite GP with first term as ‘ a ’ and common ratio as ‘ r ’, where r 1 , is,
S a (1 r ) [This point was asked in MN-2014 (2 marks)]
If a , b , c , d , are in GP then ak , bk , ck , dk , and a k , b k , c k , d k , are also in GP.
If a1 , a2 , a3 , and b1 , b2 , b3 , are two GPs then the sequence a1b1 , a2 b2 , a3b3 , is also in GP.
Three terms in a GP are taken as a r , a, ar . Four terms in a GP are taken as a r 3 , a r , ar , ar 3 .
Five terms in a GP are taken as a r 2 , a r , a, ar , ar 2 .
If every term of a GP is raised by the same index then the new progression is also in GP.
The product of the terms which are equidistance from the beginning and from the end is constant
is equal to the product of 1st and last term.
If a , b , c , d are in GP then (b a ) (c b) ( d c ) b 2 ac, c 2 bd , ad bc
If every term of a GP is multiplied or divided by some non-zero fixed quantity, the resulting
progression is also in GP.
1 1 1 12
S 7 3 2 3 7 3 10
2 2 2 1 1 2
Example 0.22 [EE-2016 (2 marks)]: Let S n 0 n n , where 1 . The value of in the range
0 1 , such that S 2 is _____.
2 3 4
Solution: S n 0 n n 0 2 2 3 3 4 4 S 2 3 4 …(i)
S 2 2 3 3 4 4 5 …(ii)
(i) (ii) (1 ) S 2 3 4 5 (1 ) S (1 ) 2
2 2 4 ( 4) 2 4(2)(1)
Now, as S 2 , so 2 2
2(1 ) 1 2 4 1 0
(1 ) 2(2)
1.707, 0.292 . As (0,1) , so 0.292 .
Harmonic Progression: A set of terms is said to be in harmonic progression when their reciprocals
1 1 1
are in arithmetic progression. The general form of a harmonic progression is , , , ,
a a d a 2d
where a and d are the first term and common difference, respectively, of arithmetic progression.
The nth term of a HP is, tn 1 {a ( n 1) d }
Three terms in a HP can be taken as 1 ( a d ) ,1 a ,1 ( a d )
Four terms in a HP can be taken as 1 ( a 3d ) ,1 ( a d ) ,1 ( a d ) ,1 ( a 3d )
Five terms in a HP can be taken as 1 ( a 2 d ) ,1 ( a d ) ,1 a ,1 ( a d ) ,1 (a 2d )
If a1 , a2 , a3 , ( ai 0i ) are in GP then log x a1 , log x a2 , log x a3 , are in AP and vice-versa;
also log a1 x, log a2 x, log a3 x, are in HP and vice-versa
Means: Let A be the arithmetic mean, G be the geometric mean and H be the Harmonic mean
between two positive quantities a and b , then
ab 2ab 2
A , G ab , H A, G , H are in GP, i.e., G AH AG H
2 ab
If we have n quantities a1 , a2 , , an , then
Arithmetic mean of n quantities is, AM ( a1 a2 an ) n
Geometric mean of n quantities is, GM (a1a2 an )1 n
1 1 1
Harmonic mean of n quantities is, HM n
a1 an a2
AM GM HM ( a1 a2 an ) n ( a1a2 an )1 n and equality holds when
a1 a2 an .
Natural Numbers
Sum of first n natural numbers is, 1 2 n n (n 1) 2
Sum of squares of first n natural numbers is, 12 22 n 2 n( n 1)(2n 1) 6
Sum of cube of first n natural numbers is, 13 23 n3 n 2 ( n 1) 2 4
Example 0.23 [CS-2008 (2 marks)]: Let P i and Q i , where K (a) P QK
1 i 2 K 1 i 2 K (b) P QK
i odd i even
is a positive integer. Then (c) PQ
(d) P Q 2K
K
Solution (a): P i 1 3 5 (2 K 1) {2 ( K 1)2} K 2
1 i 2 K 2
i odd
K
Q i 2 4 6 2k {4 ( K 1)2} K 2 K . Thus Q P K P Q K .
1 i 2 K 2
i even
Please note that the questions on ‘Sequence, Series and Progression’ also came in previous years
‘General Aptitude’ section in GATE examination. These questions are given in ‘General
Aptitude’ section.
Exercise: 0.2
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. Let Tr be r th term of an A.P. whose first term is a and common difference is d . If for some
positive integers m , n , m n , Tm 1 n and Tn 1 m , then a d equals
1 1 1
(a) (b) 1 (c) (d) 0
m n mn
2. The 19th term from the end of the series 2 6 10 86 is _____.
3. In a certain A.P., 5 times the 5th term is equal to 8 times the 8th term, then its 13th term is
_____.
4. If 7th and 13th term of an A.P. be 34 and 64 respectively, then its 18th term is _____.
5. The nth term of the series 3 10 17 and 63 65 67 are equal, then the value of n
is _____.
6. 7th term of an A.P. is 40, then the sum of first 13 terms is _____
7. The first term of an A.P. is 2 and common difference is 4. The sum of its 40 terms will be
_____.
8. The sum of the first and third term of an A.P. is 12 and the product of first and second term is
24, the first term is _____
9. S S r 1
If Sr denotes the sum of the first r terms of an A.P., then 3 r is equal to
S 2 r S 2 r 1
(a) 2r 1 (b) 2r 1 (c) 4r 1 (d) 2 r 3
10. If the sum of the first 2n terms of 2, 5,8, is equal to the sum of the first n terms of
57,59, 61, , then n is equal to _____
11. If the sum of the 10 terms of an A.P. is 4 times to the sum of its 5 terms, then the ratio of first
term and common difference is
(a) 1 : 2 (b) 2 : 1 (c) 2 : 3 (d) 3 : 2
12. 150 workers were engaged to finish a piece of work in a certain number of days. 4 workers
dropped the second day, 4 more workers dropped the third day and so on. It takes eight more
days to finish the work now. The number of days in which the work was completed is _____.
13. If the sum of first p terms, first q terms and first r terms of an A.P. be x , y and z
x z y
respectively, then (q r )
( r p) ( p q ) is
p q r
(a) 1 (b) pqr (c) xyz (d) 0
14. The sum of all odd numbers of two digits is _____.
15. If sum of n terms of an A.P. is 3n 2 5n and Tm 164 , then m _____
16. 1 1 1
The sum of n terms of the series is
1 3 3 5 5 7
1 1
(a) 2n 1 (b) 2n 1 (c) 2n 1 (d) ( 2n 1 1)
2 2
17. After inserting n A.M.’s between 2 and 38, the sum of the resulting progression is 200. The
value of n is _____
18. If a, b, c, d , e, f are six A.M.’s between 2 and 12, then a b c d e f is equal to _____.
19. If a1 , a2 , a3 , , a24 are in arithmetic progression and a1 a5 a10 a15 a20 a24 225 , then
a1 a2 a3 a23 a24 _____.
20. If a, b, c are in A.P., then 1 (bc) ,1 (ca ) ,1 ( ab ) will be in
(a) AP (b) GP (c) HP (d) None of these
21. n n
If log 2, log(2 1), log(2 3) are in A.P., then n _____
(a) 5 2 (b) log 2 5 (c) log 3 5 (d) 3 4
22. The numbers ( 2 1),1, ( 2 1) will be in
(a) AP (b) GP (c) HP (e) None of these
23. th th th
If the p , q and r term of a G.P. are a, b, c respectively, then a b r p c p q is equal to
qr
37. 3 5 7
1 2
is equal to _____.
2 2 23
38. Sum of the series 1 2.2 3.2 2 4.23 100.299 is
(a) 100.2100 1 (b) 99.2100 1 (c) 99.2100 1 (d) 100.2100 1
39. The sum of the series 3 33 333 n terms is
(10n 1 9n 28) (10 n 1 9n 10) (10 n 1 10n 9)
(a) (b) (c) (d) None of these
27 27 27
40. The sum of n terms of the following series 1 (1 x ) (1 x x 2 ) will be
1 xn x(1 x n ) n (1 x) x (1 x n )
(a) (b) (c) 2
(d) None of these
1 x 1 x (1 x )
41. The sum to n terms of the series 1 3 7 15 31 is
(a) 2n 1 n (b) 2n 1 n 2 (c) 2n n 2 (d) None of these
42. The sum of 13 23 33 43 153 is
(a) 22000 (b) 10000 (c) 14400 (d) 15000
43. th
A series whose n term is ( n x) y , the sum of r terms will be
(a) r ( r 1)
2x ry (b) r ( r 1)
2x (c) r ( r 1)
2x ry (d) r ( r 1)
2x rx
( p 1) th term from the end is ( n p 1)th term from the beginning in the expansion of ( x y ) n .
n
Cr n r 1
n
Cr n Cn r {n !} [{(n r )!}{r !}] ; n Cr n C r 1 n 1Cr ; r n C r n n 1Cr 1 ; n
.
Cr 1 r
n n
If Cx C y , then either x y or x y n
n
(1 x ) n n C0 x 0 n C1 x1 n C2 x 2 n C3 x 3 n Cn x n r 0 n Cr x r , n N
n n
Coefficient of x n r y r in ( ax by ) n r 0 n Cr ( ax) n r (by ) r r 0 n Cr a n r b r x n r y r is
n
Cr a n r b r .
n! P P P
( x1 x2 x3 xr ) n x1 1 x2 2 xr r , where P1 P2 Pr n and
P1 ! P2 ! Pr !
0 P1 , P2 , , Pr n .
n r 1
The total number of terms in ( x1 x2 x3 xr ) n is Cr 1 .
n
In the expansion of (1 x)
If x 1 , then n C0 n C1 n C2 n C3 n Cn 2n
If x 1 , then n C0 n C1 n C2 n C3 ( 1) n n Cn 0
n
C0 n C2 n C4 n C1 n C3 n C5 2n 1
n
The middle term in Binomial Expansion of ( x y ) n r 0 n Cr x n r y r , n N , depends on the
value of n
If n is even, then the total number of terms in the expansion is odd. So, there is only one
middle term, i.e. {( n 2) 1}th term is the middle term
If n is odd, then the total number of terms in the expansion is even. So, there are two middle
terms, i.e. {( n 1) 2}th and {( n 3) 2}th term are the two middle terms.
Binomial Theorem for any index: Let n be a rational number and x ( | x | 1 ) be a real number.
n(n 1) n( n 1)(n 2) n(n 1)(n 2) ( n r 1)
Then (1 x ) n 1 nx x2 x3
xr .
2! 3! r!
The condition | x | 1 is unnecessary if n is a whle number, while the same condition is
essential if n is a rational number other than a whole number.
There are infinite number of terms in the expansion of (1 x) n , when n is negative integer or
a fraction.
(1 x ) 1 1 x x 2 x 3 ( 1) r x r
(1 x ) 1 1 x x 2 x 3 x r
(1 x) 2 1 2 x 3x 2 4 x 3
(1 x ) 2 1 2 x 3x 2 4 x 3
Exercise: 0.3
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
0.4 Trigonometry
Angle in Degrees: When two lines intersect, four angles are formed; if all four angles are equal,
then each angle is called a right angle. When a right angle is divided into 90 equal parts, each part is
called a degree. A degree is divided into 60 minutes. A minute is divided into 60 seconds; i.e., 1 right
angle 90 degrees 90 o ; 1 degree 60 minutes 60 ' ; 1 minute 60 seconds 60 ''
Angle in Radians: The angle subtended at the centre of any circle by an arc of that circle is the ratio
length of arc
of length of an arc to the radius of the circle, i.e. (in radians) . Hence, one radian is
radius of circle
the angle subtended at the centre of any circle by an arc equal in length to the radius of that circle.
Conversion of degree into radians: 1 degree ( 180) radians
Conversion of radians into degree: 1 radian (180 ) degrees
Angle measured in anti-clockwise is taken as ve ; Angle measured in clockwise is taken as ve .
Trigonometric function of an angle: The six trigonometric ratios sine, cosine, tangent, cotangent,
secant and cosecant of an angle , 0o 90 o are defined as the ratios of two sides of a right-
angled triangle with as
one of the angles. The
opposite side of the angle
, for which we find the
trigonometric ratios, is
taken as ‘perpendicular’;
the side which makes
right angle with the
perpendicular is taken as
‘base’; and the remaining Figure 0.4: (a) The Right angled triangle (b) Sign of Trigonometric Functions
third side is called
Perpendicular BC 1
hypotenuse, as shown in Fig. 0.4; then by definition, sin ;
Hypotenuse AC cosec
Base AB 1 Perpendicular BC 1
cos ; tan .
Hypotenuse AC sec Base AB cot
(a) (b)
(c) (d)
(e) (f)
The table for the formulae for Allied angles is given as:
o o o o o o
180 180 90 90 270 270
sin sin sin sin cos cos cos cos
cos cos cos cos sin sin sin sin
tan tan tan tan cot cot cot cot
cot cot cot cot tan tan tan tan
sec sec sec sec cosec cosec cosec cosec
cosec cosec cosec cosec sec sec sec sec
The results in the above table are not to be memorised. We have to follow the following procedure:
Step 1: First find the quadrant in which the angle lies and get the sign of the trigonometric function in
o
that quadrant. For e.g., for sin(270 ) , assume is very small, 270o lie in the 4th quadrant
and sin is negative in 4th quadrant.
Step 2: Now divide the given angle by 90 o and by dividing we get and even or odd number. If we get
even number then the trigonometric function remains unchanged. If we get an odd number then the
trigonometric functions interchange with each other according as: sin cos , tan cot ,
o
sec cos ec . For e.g., for sin(270 ) , if we divide 270o by 90 o , we get 3 which is an odd
o
number, so sin changes into cos . So sin(270 ) cos .
Identities relating to the angles of a triangle: In any ABC , the following results hold true:
cot B cot C cot C cot A cot A cot B 1
sin 2 A sin 2B sin 2C 4sin A sin B sin C tan A tan B tan C (tan A)(tan B)(tan C )
A B B C C A A B C
tan tan tan tan tan tan 1 sin A sin B sin C 4 cos cos cos
2 2 2 2 2 2 2 2 2
A B C A B C A B C
cos A cos B cos C 1 4 sin sin sin cot cot cot cot cot cot
2 2 2 2 2 2 2 2 2
Relations between the sides and angles of a triangle: In the following relations, A, B, C are the
angles and a, b, c are the sides of the triangle opposite to the angles A, B, C , respectively.
a b c
2 R (Sine formula), where R is the circumradius
sin A sin B sin C
2 2 2 2 2 2 2 2 2
c a b 2ab cos C ; b c a 2ac cos B ; a b c 2bc cos A (cosine formula)
a c cos B b cos C ; b a cos C c cos A ; c b cos A a cos B (Projection formula)
Half angle formula: In the following formula A, B, C are the angles of a triangle and a, b, c are the
sides of the triangle opposite to the angles A, B, C , respectively. Also, perimeter of the triangle
2s a b c .
A ( s b)( s c ) B ( s c )( s a ) C ( s a )( s b)
sin ; sin ; sin
2 bc 2 ca 2 ab
A s ( s a) B s ( s b) C s (s c)
cos ; cos ; cos
2 bc 2 ca 2 ab
A ( s b)( s c ) B ( s c )( s a ) C ( s a )( s b)
tan ; tan ; tan
2 s( s a ) 2 s( s b) 2 s(s c)
B C bc A CA ca B AB a b C
tan cot ; tan cot ; tan cot (Napier’s formula)
2 bc 2 2 ca 2 2 ab 2
Inradius and Exradii of a triangle: Let a, b, c are the sides of the triangle opposite to the angles
A, B, C , respectively. Let r be the inradius of the triangle and r1 , r2 , r3 are the exradius of the
triangle. An escribed circle is a circle which touches one of the sides of a triangle internally and the
other two sides externally (sides produced). So r1 , r2 , r3 are the exradius of the triangle ABC opposite
to angles A, B, C , respectively. In all the following formulae, is the area of triangle ABC , R is the
circumcentre of the triangle ABC .
A B C A B C
r ( s a ) tan ( s b) tan ( s c) tan 4 R sin sin sin
s 2 2 2 2 2 2
A A B C B A B C
r1 s tan 4 R sin cos cos r2 s tan 4 R cos sin cos
sa 2 2 2 2 s b 2 2 2 2
C A B C
r3 s tan 4 R cos cos sin r1 r2 r3 r 4R
sc 2 2 2 2
Example 0.25 [CS-1995 (1 mark)]: In the interval [0, ] the equation x cos x has
(a) No solution (b) Exactly one solution
(c) Exactly two solutions (d) An infinite number of solutions
Solution (b): By sketching the graph of y x
and y cos x , we have only one point of
intersection. So the equation x cos x has
exactly one solution.
Inverse Circular Function: The inverse of a function f : A B exists if f is one-one onto i.e.,
1
a bijection and is given by f ( x) y f ( y) x . Consider the sine function with domain and
range [ 1,1] . Clearly this function is not a bijection and so it is not invertible. If we restrict the
domain of it in such a way that it becomes one–one, then it would become invertible. If we consider
sine as a function with domain [ 2 , 2] and co-domain [ 1,1] , then it is a bijection and
therefore, invertible. The inverse of all trigonometric function along with domain, range, nature, and
their graphs is given in the following table:
1 x, 0 x 1 x, x 0, x 2
cos (cos x ) sec (sec x)
2 x, x 2 x, 0 x , x 2
tan(tan 1 x) x for all x 1
cosec(cosec x) x for all x 1
x, 0 x 2 x, 0 x 2
1 1
tan (tan x) x , 2 x 3 2 cosec (cosec x ) x, x ( 2 , 3 2]
x 2 , 3 2 x 2 x 2 , 3 2 x 2
1 1 1 1
sin (sin x) , cos (cos x) , sec (sec x) and cosec (cosec x) are periodic with period 2 .
1 1
While, tan (tan x) and cot (cot x) are periodic with period .
x y x y x y x y
sinh x sinh y 2 sinh cosh sinh x sinh y 2 cosh sinh
2 2 2 2
x y x y x y x y
cosh x cosh y 2 cosh cosh cosh x cosh y 2 sinh sinh
2 2 2 2
x
2 sinh x cosh y sinh ( x y ) sinh ( x y ) cosh x sinh x e
2 cosh x cosh y cosh ( x y ) cosh ( x y ) cos h x sin h x e x
2 cos h x sin h y sin h ( x y ) sin h ( x y) 2cos h2 x cos h 2 x 1
2 sinh x sinh y cosh ( x y ) cosh ( x y ) 2 tanh x
n
sinh 2 x 2 sinh x cosh x
(cos h x sin h x) cos h nx sin h nx 1 tanh 2 x
2 2 2 2 1 tan h 2 x
cosh 2 x cosh x sinh x 2 cosh x 1 1 2 sinh x 2
1 tan h x
3 tanh x tanh 3 x 2 tan h x
tanh 3 x tan h 2 x 2
1 3 tanh 2 x 1 tan h x
sinh 3x 3sinh x 4sinh 3 x cosh 3x 4 cosh3 x 3cosh x
cos h x sin h x e x n
(cosh x sinh x) cosh nx sinh nx
sinh( x) sinh x cosh( x ) cosh x tanh( x ) tanh x
coth( x ) coth x sec h( x ) sec hx cosec h( x ) cosec hx
2 2
Transformation of Hyperbolic Function: Since cosh x sinh x 1 sinh x cosh 2 x 1
2
1 sech x tanh x 1
sinh x sinh x sinh x
. In a similar manner we can
sec h x 2 2
1 tanh x coth x 1
express cosh x , tanh x , coth x , etc., in terms of other hyperbolic functions.
Relation between Hyperbolic and Circular Function: We have from Euler formulae, (1)
ix ix ix ix
e cos x i sin x (2) e cos x i sin x . Adding (1) and (2) cos x (e e ) 2 . Subtracting
(2) from (1) sin x (eix e ix ) 2i . Replacing x by ix in these values, we get
x
cos(ix ) (e x
e ) 2 cosh x cos(ix) cosh x , sin(ix ) (e e x ) (2i ) i (e x e x ) 2
x
sin ix i sinh x
sin(ix) i sinh x . Also tan ix i tanh x . Similarly replacing x by ix in the
cos ix cosh x
definitions of sinh x and cosh x , we get cos h (ix) (eix e ix ) 2 cos x . Thus, we obtain the
following relations between hyperbolic and trigonometrical functions: sin(ix) i sinh x ;
sinh(ix ) i sin x ; sinh x i sin(ix) ; sin x i sinh(ix) ; cos(ix) cosh x ; cosh(ix ) cos x ;
cosh x cos(ix ) ; cos x cosh(ix ) ; tan(ix ) i tanh x ; tanh(ix ) i tan x ; tanh x i tan(ix) ;
tan x i tanh(ix) ; cot(ix ) i coth x ; coth(ix) i cot x ; coth x i cot(ix ) ; cot x i coth(ix ) ;
sec(ix ) sec h x ; sec h (ix) sec x ; sec h x sec(ix ) ; sec x sec h(ix ) ; cosec(ix ) i cosec h x ;
cos ech (ix ) i cosec x ; cosec h x i cosec (ix) ; cosec x i cosec h(ix ) .
e x e x x3 2 5 17 7
tanh x x x
x x x , x
e e 3 15 315 2
3
1 x x 2 5
coth x x , 0 x
x 3 45 945
x 2 5 4 61 6
sec hx 1 x x , x
2 24 720 2
3
1 x 7x
cosec hx , 0 x
x 6 360
Example 0.26 [EC-2007 (1 mark)]: For x 1 , coth( x) can be approximated as
2
(a) x (b) x 2 (c) 1 x (d) 1 x
1 x x3 2 5
Solution (c): From the expansion of coth x x , except the first term all
x 3 45 945
other terms tends to zero when x 1 , so For x 1 , coth( x) can be approximated as 1 x .
Period of Hyperbolic Functions: If for any function f ( x) , f ( x T ) f ( x ) , then f ( x) is called
the Periodic function and least positive value of T is called the Period of the function.
sinh x sinh(2 i x) , cosh x cosh(2 i x ) and tanh x tanh( i x ) . Therefore the period of
these functions are respectively 2 i , 2 i and i . Also period of cos ech x , sec h x and coth x are
respectively 2 i , 2 i and i .
Remember that if the period of f ( x ) is T , then period of f ( nx ) will be T n .
Hyperbolic function are neither periodic functions nor their curves are periodic but they show the
algebraic properties of periodic functions and having imaginary period.
Exercise 0.4
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. The circular wire of radius 7 cm is cut and bend again into an arc of a circle of radius 12 cm.
The angle (in degrees) subtended by an arc at the centre of the circle is _____.
2. The equation sec 2 4 xy ( x y ) 2 is only possible when
(a) x y (b) x y (c) x y (d) None of these
3.
For 0 , if x cos 2 n , y sin 2 n , z (cos 2 n )(sin 2 n ) , then
2 n 0 n 0 n 0
(a) xyz xz y (b) xyz xy z (c) xyz x y z (d) Both (b) and (c)
4. Which of the following number is rational?
(a) sin 15 o (b) cos15 o (c) sin15o cos15o (d) sin15 o cos 75o
5. If tan m ( m 1) and tan 1 (2 m 1) then _____ degrees.
6. The value of tan x tan 3 x whenever defined never lie between
(a) 1 3 and 3 (b) 1 4 and 4 (c) 1 5 and 5 (d) 5 and 6
7. If A cos 2 sin 4 , then for all values of
(a) 1 A 2 (b) 13 16 A 1 (c) 3 4 A 13 16 (d) 3 4 A 1
8. 2
The number of solutions of the equation 3sin x 7 sin x 2 0 in the interval [0, 5 ] is ___
9. Number of solutions of the equation tan x sec x 2 cos x , lying in the interval [0, 2 ] is ___
10. General value of satisfying the equation tan 2 sec 2 is
(a) m , n ( 3) (b) m , n ( 3) (c) m , n ( 6) (d) None of these
11. The number of integral values of k , for which the equation 7 cos x 5 sin x 2 k 1 has a
solution is _____.
12. (a) 3 : (2 3) (b) 1: 6
If the angles of a triangle are in the ratio 4 :1 :1 , then the
ratio of the longest side to the perimeter is (c) 1: (2 3) (d) 2 : 3
13. In a triangle ABC , B 3 and B 4 and D divides BC internally in the ratio 1: 3 .
Then sin BAD sin CAD is equal to
(a) 1 3 (b) 1 3 (c) 1 6 (d) 2 3
14. In a ABC , 2ac sin ( A B C ) 2 is equal to
(a) a 2 b 2 c 2 (b) c 2 a 2 b 2 (c) b 2 c 2 a 2 (d) c 2 a 2 b 2
15. cos A cos B cos C
In a ABC , and a 2 then square of the area of a triangle is _____
a b c
16. In a triangle ABC , let C 2 . If r is the in-radius and R is the circum-radius of the
triangle, then 2( r R) is equal to
(a) a b (b) b c (c) c a (d) a b c
17. A regular polygon of nine sides, each of length 2 is inscribed in a circle. The radius of the
circle is
(a) cosec( 9) (b) cosec( 3) (c) cot( 9) (d) tan( 9)
18. From the top of a light house 60 m high with its base at the sea level the angle of depression of
a boat is 15o. The distance (in metres) of the boat from the foot of the light house is
3 1 3 1 3 1
(a) 60 (b) 60 (c) (d) None of these
3 1 3 1 3 1
19. A person observes the angle of deviation of a building as 30 o . The person proceeds towards
the building with a speed of 25( 3 1) m/h. After 2 hours, he observes the angle of elevation
as 45o . The height of the building (in metres) is
20. The shadow of a tower standing on a level ground is found to be 60 m longer when the sun’s
altitude is 30 o than when it is 45o . The height (in metres) of the tower is
(a) 60 m (b) 30 m (c) 60 3 (d) 30( 3 1)
21. A person is standing on a tower of height 15( 3 1) m and observing a car coming towards
the tower. He observed that angle of depression changes from 30 o to 45o in 3 sec. What is
the speed (in km/h) of the car? _____
22. The angle of elevation of the top of a pillar at any point A on the ground is 15o. On walking
40 m towards the pillar, the angle becomes 30o. The height (in metres) of the pillar is _____.
23. A man from the top of a 100 metre high tower looks a car moving towards the tower at an angle of
depression of 30o. After some time, the angle of depression becomes 60o. The distance (in metre)
travelled by the car during this time is
(a) 100 3 (b) 200 3 (c) 100 3 (d) 200 3
24. For a man, the angle of elevation of the highest point of the temple situated east of him is 60o .
On walking 240 metres to north, the angle of elevation is reduced to 30 o then the height of the
temple is
(a) 60 6 (b) 60 (c) 50 3 (d) 30 6
25. 2
Principal value of sin 1 sin is (a) 2 3 (b) 2 3 (c) 4 3 (d) 3
3
26. Considering only the principal values, if tan(cos 1 x ) sin cot 1 (1 2) , then x is equal to
Centroid of a triangle: The centroid of a triangle is the point of intersection of its medians, as
shown in Fig. 0.6. The centroid divides the medians in the ratio 2 : 1 (Vertex : base). If A( x1 , y1 )
, B( x2 , y2 ) and C ( x3 , y3 ) are the vertices of a triangle. If G be the centroid upon one of the median
x1 x2 x3 y1 y2 y3
(say) AD, then AG : GD 2 :1 Co-ordinate of G are , .
3 3
Incentre: The incentre of a triangle is the point of intersection of internal bisector of the angles,
as shown in Fig. 0.8. Also it is a centre of a circle touching all the sides of a triangle. Co-ordinates
ax bx2 cx3 ay1 by2 cy3
of incentre 1 , , where a, b, c are the sides of triangle ABC .
abc abc
Excircle: A circle touches one side outside the triangle and other two extended sides then circle is
known as excircle, as shown in Fig. 0.9. Let ABC be a triangle then there are three excircles with
three excentres. Let I1 , I 2 , I 3 opposite to vertices A( x1 , y1 ) , B( x2 , y2 ) and C ( x3 , y3 ) then
ax1 bx2 cx3 ay1 by2 cy3
coordinates of three excentres can be given as: I1 , ;
a b c a b c
ax bx2 cx3 ay1 by2 cy3 ax1 bx2 cx3 ay1 by2 cy3
I2 1 , ; I3 , .
a b c abc abc abc
BD AB c
Angle bisector divides the opposite sides in the ratio of remaining sides, i.e.
DC AC b
Incentre divides the angle bisectors in the ratio (b c) : a, (c a) : b and (a b ) : c
Figure 0.8: Incentre of a triangle Figure 0.9: Excircle of a triangle Figure 0.10: Orthocentre of a triangle
If a triangle is right angled triangle, then orthocentre is the point where right angle is formed.
If the triangle is equilateral then centroid, incentre, orthocentre, circum-centre coincides.
Orthocentre, centroid and circum-centre are always collinear and centroid divides the line
joining orthocentre and circum-centre in the ratio 2 : 1
In an isosceles triangle centroid, orthocentre, incentre, circum-centre lie on the same line.
Example 0.28 [MN-2010 (2 Solution (a): The given point are collinear if the area of triangle
marks)]: The value of k for 5 5 1
which the points (5, 5) , ( k ,1) , 1
formed by given points is zero, i.e. k 1 1 0
(10, 7) lie on a straight line is 2
(a) 5 (b) 5 10 7 1
(c) 2 (d) 2 5(1 7) 5( k 10) 1(7 k 10) 0 k 5 .
Rotation of axes without changing the origin: Let O be the origin. Let P ( x , y ) with respect
to axes OX and OY and let P ( x, y ) with respect to axes OX and OY are formed by rotating
the axis without shifting of origin as shown in Fig. 0.12(a), where X OX YOY then
following relations holds true: x x cos y sin , y x sin y cos ; and
x x cos y sin , y x sin y cos .
The above relation between ( x, y ) and ( x, y ) x y
can be easily obtained with the help of x cos sin
following table y sin cos
Figure 0.12: (a) Rotation of Axis without Shifting of Origin. (b) Rotation of Axis and Shifting of Origin
Change of origin and rotation of axes: If origin is changed to O ( , ) and axes are rotated
about the new origin O by an angle in the anticlock-wise sense, as shown in Fig. 0.12(b), such
that the new co-ordinates of P ( x, y ) become ( x, y ) then the equations of transformation will be
x x cos y sin and y x sin y cos .
Reflection (Image of a point): Let ( x, y ) be any point, then its image with respect to
x axis is ( x, y ) y axis is ( x, y )
Origin is ( x, y ) Line y x is ( y, x )
Locus: When a point moves in accordance with a geometric law, its path is called a locus. The
equation to the locus is the equation connecting the x and y coordinates of every point on the curve.
x y
Intercept form: The line making an intercept a and b on the x and y axis is 1
a b
Normal form: If the perpendicular from the origin to the line is of length p and makes an angle
with the positive direction of the x axis, the equation of the line is x cos y sin p
Line parallel to the x axis is y constant; line parallel to the y axis is x constant
Line passes through the origin having slope m is y mx
Line passes through the point ( x1 , y1 ) with slope m is y y1 m( x x1 )
y 2 y1 y2 y1
Line passes through the points ( x1 , y1 ) and ( x2 , y2 ) is y y1 ( x x1 ) , where is
x2 x1 x2 x1
the slope of the line
The line passing thought a fixed point ( x1 , y1 ) and having inclination with the positive
x x1 y y1
direction of x axis is r , where r is the distance of any point ( x , y ) , on the
cos sin
line, from the given point ( x1 , y1 ) . The coordinates of any point ( x , y ) on this line are given by
x x1 r cos , y y1 r sin
Position of a point with respect to a line: Two points ( x1 , y1 ) and ( x2 , y2 ) lies on the
ax1 by1 c
same side of the straight line ax by c 0 according as 0
ax2 by2 c
ax1 by1 c
opposite sides of the straight line ax by c 0 according as 0
ax2 by2 c
Two or more lines: The coordinates of the point of intersection of two straight lines are obtained by
solving the equations of the two lines
If is the acute angle between the lines y m1 x c1 and y m2 x c2 , then
m tan ( m2 m1 ) (1 m1m2 ) .
If the two lines are parallel then the angle between them is 0 o m1 m2
If the two lines are perpendicular then the angle between them is 90 o 1 m1m2 0
The equation of a line parallel to ax by c 0 is ax by k 0 , where k is found by given
condition.
The equation of a line parallel to ax by c 0 is bx ay k 0 , where k is found by given
condition.
The two lines a1 x b1 y c1 0 and a2 x b2 y c2 0 are
Coincident if a1 a2 b1 b2 c1 c2 Parallel if a1 a2 b1 b2 c1 c2
Intersecting if a1 a2 b1 b2 Perpendicular is a1a2 b1b2 0
The equation of any line passing through the point of intersection of lines a1 x b1 y c1 0 and
a2 x b2 y c2 0 is given as a1 x b1 y c1 ( a2 x b2 y c2 ) 0 , where is a constant and
found by applying the given condition.
Example 0.29 [MN-2007 (1 mark)]: If the slope of a diagonal of a rectangle is m , the slope of the
other diagonal is
(a) 1 2m (b) 1 2m (c) 1 m (d) 1 m
Solution: Let the coordinated of any rectangle ABCD are A( a, b) , B (c, b) ,
C (c, d ) , D( a, d ) , as shown in figure. The slope of diagonal AC is
d b d b
mAC m (say) and slope of diagonal BD is mBD m ; thus
ca ac
given options are wrong.
Example 0.30 [AG-2012 (1 mark)]: The line y x 1 can be expressed in polar coordinates ( r , )
as
(a) r cos (b) r sin (c) r (cos sin ) 1 (d) r (cos sin ) 1
Solution (d): In polar coordinates, we have x r cos and y r sin . So the line y x 1 in polar
coordinates will be r sin r cos 1 r (cos sin ) 1 .
Example 0.31 [IN-2016 (1 mark)]: A straight line of the form y mx c passes through the origin
and the point ( x, y ) (2, 6) . The value of m is _____.
Solution: As the given straight line passes through the origin, so c 0 , thus we have y mx . So
y mx passes through (2, 6) , thus we have 6 m(2) m 3 .
Pair of Straight Lines: Let the equation of two lines be ax by c 0 …(i) and
ax by c 0 …(ii). Hence ( ax b y c)( ax by c) 0 is called the joint equation of lines
(i) and (ii) and conversely, if joint equation of two lines be ( ax by c)( ax by c) 0 then
their separate equation will be ax by c 0 and ax by c 0 .
Equation of a pair of straight lines passing through origin: The equation ax 2 2hxy by 2 0
represents a pair of straight line passing through the origin where a, h, b are constants. Let the
lines represented by ax 2 2hxy by 2 0 be y m1 x 0 and y m2 x 0 , where,
h h 2 ab h h 2 ab 2h a
m1 and m2 then, m1 m2 and m1m2 . Then, two
b b b b
straight lines represented by ax 2 2hxy by 2 0 are ax hy y h 2 ab 0 = 0.
The lines are real and distinct if h2 ab 0
The lines are real and coincident if h2 ab 0
The lines are imaginary if h2 ab 0
If the pair of straight lines ax 2 2hxy by 2 0 and ax 2 2hxy by 2 0 are common, then
( ab ab) 2 4( ah ah)( hb hb) .
The equation of the pair of straight lines passing through origin and perpendicular to the pair
of straight lines represented by ax 2 2hxy by 2 0 is given by bx 2 2hxy ay 2 0 .
Angle between the pair of lines: The angle between the lines represented by
2 2
2 h ab 2 h ab
ax 2 2hxy by 2 2 gx 2 fy c 0 is given by tan tan 1 .
ab ab
The pair of lines ax 2 2hxy by 2 2 gx 2 fy c 0 and ax 2 2hxy by 2 0 are parallel to
each other thus the angle between the pair of lines represented by ax 2 2hxy by 2 0 is
given by tan 2 h 2 ab ( a b) .
The lines are parallel if the angle between them is zero, i.e. 0 tan 0 h 2 ab
The lines are perpendicular if the angle between them is 2 , i.e.
2 tan 1 0 a b 0
Point of intersection of pair of lines: The point of intersection of pair of line represented by
ax 2 2hxy by 2 2 gx 2 fy c 0 is found by the following steps:
Step I: Let ax 2 2hxy by 2 2 gx 2 fy c 0
Step II: Find x 2 ax 2 hy 2 g 0 (by keeping y as constant)
Step III: Find y 2hx 2by 2 f 0 (by keeping x as constant)
Example 0.32 [CS-1995 (1 mark)]: The value of k for which 4 x 2 8 xy ky 2 0 does not represent
a pair of straight lines (both passing through the origin) is
(a) 0 (b) 2 (c) 9 (d) 3
Solution: On representing the given equation in the form of quadratic equation in x , we have
8 y 64 y 2 4(4)(ky 2 ) 8 y 64 y 2 16ky 2 8y 4y 4 k
4 x 2 8 xy ky 2 0 x
. Now
2(4) 8 8
the given equation will represent the pair of straight lines if 4 k 0 k 4 . Thus from given
options, if k 9 4 then the given equations will not represent pair of straight lines.
Circle: The equation to the circle having its centre at O ( h, k ) and radius r , as shown in
Fig. 0.13, is ( x h) 2 ( y k ) 2 r 2 . Any point P on this circle is ( h r cos , k r sin )
, where is the angle made by OP with the x axis. Figure 0.13:
If the centre is the origin and radius is r , the equation of the circle is x 2 y 2 r 2 . A circle
Any point on this circle can be taken as ( r cos , r sin ) , where is the parameter of the point.
The general equation to the circle is x 2 y 2 2 gx 2 fy c 0 . Its centre is ( g , f ) and radius
is g2 f 2 c .
Equation to the circle described on the line joining the points ( x1 , y1 ) and ( x2 , y2 ) as a diameter
is ( x x1 )( x x2 ) ( y y1 )( y y2 ) 0
The second degree curve in two variables ax 2 2hxy by 2 2 gx 2 fy c 0 represents a circle
if a b and h 0 . To find the centre and radius of this circle, reduce it to the form
2g 2f c g f g 2 f 2 ca
x2 y 2 x y 0 centre , and radius is
a a a a a a
A circle can be made to satisfy any three geometrical conditions such as
passing through three given points or touching three given lines
passing through two given points and touching a given line
passing through a given point and touching two given lines
passing through two given points and having its centre on a given line
The length of tangent from point P ( x1 , y1 ) to the circle x 2 y 2 2 gx 2 fy c 0 is
Two or more Circles: Let d be the distance between the centres of two circles with radii r1 , r2
If d r1 r2 , one circle lies completely inside the other circle
If d r1 r2 , the two circles do not intersect
If r1 r2 d r1 r2 , the two circles intersect
If d r1 r2 , the two circles touch one another externally and the point of contact divides the line
joining the centres internally in the ratio r1 : r2 .
If d r1 r2 , the two circles touch one another internally and the point of contact divides the line
joining the centres externally in the ratio r1 : r2 .
Let S x 2 y 2 2 gx 2 fy c 0 and S x 2 y 2 2 g1 x 2 f1 y c1 0 be the equations of
two circles. Then the two circles will cut each other orthogonally if 2 gg1 2 ff1 c c1
If S 0 and S 0 are the equations of two circles, then S S 0 is the equation of the radical
axis. However, if the two circles S 0 , S 0 coincides with the common chord. Therefore
S S 0 is the equation of common chord of the intersecting circles.
Equation of a circle passing through the points of intersection of the circles S 0 and S 0 can
be written as S kS 0 , where k is a real number.
The radical axis of three circles taken in pairs are concurrent and the point of concurrence is the
radical centre of the three circles.
Parabola: Standard equation of parabola is y 2 4ax . The following points hold true for y 2 4ax :
x axis is the axis of the parabola; y axis is the tangent at the vertex
Vertex is A(0, 0) ; Focus is S (a, 0) ; Directrix is x a 0
Parametric form of representation of a point P on y 2 4ax is
P ( at 2 , 2at )
At point P ( at 2 , 2at ) on y 2 4ax , slope of tangent is 1 t
Figure 0.14: A Parabola
Equation of tangent at P ( at 2 , 2at ) on y 2 4ax is x yt at 2 0
Equation of normal at P ( at 2 , 2at ) on y 2 4ax is y tx 2at at 3 0 ; slope of normal is t
If P ( at12 , 2at1 ) and Q ( at 22 , 2 at2 ) are the two point on the parabola y 2 4ax , then slope of the
chord PQ is 2 (t1 t2 ) . The equation of the chord PQ is 2 x y (t1 t2 ) 2at1t2 0
If the tangent at P ( at12 , 2at1 ) and at Q ( at 22 , 2 at2 ) are perpendicular, then t1t2 1 . In this case,
the tangents intersect on the directrix
If PQ is a chord passing through the focus, then t1t2 1
If the normal at t1 meets the parabola, again at t2 , then t2 t1 (2 t1 )
If the point P on y 2 4ax is taken in Cartesian form, then
equation of tangent at ( x1 , y1 ) is yy1 2a ( x x1 ) 0
equation of normal at ( x1 , y1 ) is xy1 2ay ( x1 y1 2ay1 ) 0
If the line y mx c touches the parabola y 2 4ax , the condition is c a m . So
y mx (a m ) is a tangent to the parabola for all values of m .
The chord of contact of tangents to the parabola y 2 4ax from the point ( x1 , y1 ) is
yy1 2a ( x x1 ) 0 , i.e., T 0
x2 y2 x2 y2
Ellipse: Standard equation of Ellipse is 1 . The following points hold true for 1:
a2 b2 a2 b2
x axis is the major axis of length 2a ; y axis is the minor axis of length 2b
b 2 a 2 (1 e 2 ) ( e 1 is the eccentricity of the
ellipse)
There are two foci, one S ( ae, 0) and the other
S ( ae, 0) , respectively, the corresponding
directrices are x a e and x a e
If P be any point on the ellipse, then (i)
SP S P 2a ; and (ii) SP S P CD 2 , where
CD is the semi-diameter parallel to the tangent Figure 0.15: An Ellipse
at P
x2 y2
Parametric form of representation of a point P on the ellipse 1 is P ( a cos , b sin ) .
a2 b2
x y
Equation of tangent at P is cos sin 1 0
a b
ax by
Equation of normal at P is a 2 b2
cos sin
If CD be a semi-diameter parallel to the tangent at P , then D is ( a sin , b cos ) and
2 2 2 2 2
CD a sin b cos
x2 y2
The locus of point of intersection of perpendicular tangents of the ellipse 1 is the circle
a2 b2
(called the director circle) x 2 y 2 a 2 b 2
x2 y2
The circle x 2 y 2 a 2 is called the auxiliary circle of the ellipse 1 , and it is the locus
a2 b2
of the foot of perpendicular from the centre of the ellipse to its tangents
x2 y2
Hyperbola: Standard equation of the Hyperbola is 1
a2 b2
x axis is the transverse axis of length 2a ; y axis is the
conjugate axis of length 2b
b 2 a 2 (e 2 1) ( e 1 is the eccentricity of the hyperbola)
There are two foci, one S ( ae, 0) and the other S ( ae, 0) ,
respectively, the corresponding directrices are x a e and
x a e Figure 0.16: A Hyperbola
If P be any point on the Hyperbola, then SP S P 2a
Parametric form of representation of a point P on the hyperbola is ( a sec , b tan )
x2 y2 x y
At P ( a sec , b tan ) on 2
2
1 , equation to the tangent at is
tan 1 sec
a b b a
x2 y2
The locus of point of intersection of perpendicular tangents of the hyperbola 2 2 1 is the
a b
2 2 2 2
circle (called the director circle) x y a b
x2 y2
The combined equation of the two asymptotes is 0
a2 b2
The angle between the asymptotes is given by 2 sec 1 (e)
x2 y2
Equation of the conjugate hyperbola is 1
a2 b2
The hyperbola is said to be rectangular if a b , the equation to the rectangular hyperbola is
x2 y2 a2
Exercise: 0.5
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. A triangle with vertices (4, 0) , ( 1, 1) and (3,5) is
(a) isosceles and right angled (b) isosceles but not right angled
(c) right angled but not isosceles (d) neither right angled nor isosceles
2. The line x y 4 divides the line joining the points ( 1,1) and (5, 7) in the ratio
(a) 2 : 1 , internally (b) 2 : 1 , externally (c) 1 : 2 , internally (d) 1 : 2 , externally
3. If the vertices of a triangle be (2,1) , (5, 2) and (3, 4) then its circumcentre is
(a) (13 2 ,9 2) (b) (13 4 ,9 4) (c) (9 4 ,13 4) (d) None of these
4. The incentre of the triangle with vertices (1, 3) , (0, 0) , and (2, 0) is
(a) (1, 3 2) (b) (2 3,1 3) (c) (2 3, 3 2) (d) (1,1 3)
5. Four points are A(6,3) , B ( 3, 5) , C (4, 2) and P ( x, y ) . Then ratio of area of PBC and
ABC is
(a) ( x y 2) 7 (b) ( x y 2) 2 (c) ( x y 2) 7 (d) None of these
6. The points (0,8 3) , (1,3) and (82, 30) are the vertices of a/an _____ triangle.
(a) equilateral (b) isosceles (c) right angled (d) none of these
7. The point (2,3) undergoes the following three transformation successively, (i) Reflection
about the line y x ; (ii) Transformation through a distance 2 units along the positive
direction of y-axis; (iii) Rotation through an angle of 45o about the origin in the anticlockwise
direction. The final coordinates of points are
(a) (1 2 , 7 2) (b) ( 1 2 , 7 2) (c) (1 2 , 7 2) (d) None of these
8. The ends of a rod of length l move on two mutually perpendicular lines. The locus of the point
on the rod which divides it in the ratio 1 : 2 is
26. x2 y2
Let P be a variable point on the ellipse
1 with foci F1 and F2 . If A is the area of
a2 b2
the triangle PF1 F2 , then the maximum value of A is
(a) 2abe (b) abe (c) abe 2 (d) None of these
27. x2 y2
Let E be the ellipse 1 and C be the circle x 2 y 2 9 . Let P and Q be the
9 4
points (1, 2) and (2,1) respectively. Then
(a) Q lies inside C but outside E (b) Q lies outside both C and E
(c) P lies inside both C and E (d) P lies inside C but outside E
28. The locus of the middle point of the intercept of the tangents drawn from an external point to
the ellipse x 2 2 y 2 2 between the coordinate axes, is
1 1 1 1 1 1 1 1
(a) 2
2
1 (b) 2
2
1 (c) 2
2
1 (d) 2
2
1
x 2y 4x 2y 2x 4y 2x y
29. x2 y2 x2 y2 1
If the foci of an ellipse 2
1 and a hyperbola coincide, then b 2 is ___.
16 b 144 81 25
30. If the circle x y a intersects the hyperbola xy c 2 in four points P ( x1 , y1 ) , Q ( x2 , y2 )
22 2
, R( x3 , y3 ) and S ( x4 , y4 ) then
4 4
(a) xi 0 (b) yi 0 (c) x1 x2 x3 x4 c 4 (d) None of these
i 1 i 1
31. If L is the line whose equation is ax by c . Let M be the reflection of L through the y
axis, and let N be the reflection of L through the x axis. Which of the following must be
true about M and N for all choices of a , b and c ?
(a) The x intercept of M and N are equal
(b) The y intercept of M and N are equal
(c) The slopes of M and N are equal are equal
(d) The slopes of M and N are equal are reciprocal
32. The centres of three circles x 2 y 2 10 x 9 0 , x2 y 2 6x 2 y 1 0 ,
x2 y 2 9x 4 y 2 0
(a) lie on the straight line x 2 y 5 (b) lie on the circle x 2 y 2 25
(c) do not lie on straight line (d) lie on circle x 2 y 2 x y 17 0
33. The parabola y 4 x 2 has vertex P . It intersects x axis at A and B . If the parabola is
translated from its initial position to a new position by moving its vertex along the line
y x 4 , so that it intersects x axis at B and C , then abscissa of C will be _____
34. An ellipse has semi-major axis of length 2 and semi-minor axis of length 1. It slides between
the coordinate axes in the first quadrant, while maintaining contact with both x axis and y
axis. The locus of centre of ellipse is
(a) x 2 y 2 3 (b) x 2 y 2 5 (c) ( x 2) 2 ( y 1) 2 5 (d) ( x 2) 2 ( y 1) 2 3
35. Let PQ : 2 x y 6 0 is a chord of the curve x 2 4 y 2 4 point R ( , ) that satisfy
2 2 1 0 ; such that area of triangle PQR is minimum; are given by:
2 1 2 1 2 1 2 1
(a) , (b) , (c) , (d) ,
5 5 5 5 5 5 5 5
Direction Ratios: Three numbers which are proportional to the direction cosines of a line are
called the direction ratio of that line. If a , b and c are three real numbers such that
l a m b n c k (say) l ak , m bk , n ck , then a , b and c are called direction ratios
of OP . Since l 2 m 2 n 2 1 (a 2 b 2 c 2 )k 2 1 k 1 a2 b2 c2 , so,
l a a 2 b 2 c 2 , m b a 2 b 2 c 2 , n c a 2 b 2 c 2 , where the signs should
be taken all positive or all negative.
Direction ratios are not unique, whereas d.c.’s are unique. i.e., a 2 b 2 c 2 1
Let r a i b j c k be a vector. Then its d.r.’s are a, b, c
r
If a vector r has d.r.’s a, b, c then r (a i b j c k )
2 2 2
a b c
D.c.’s and d.r.’s of a line joining two points: The direction ratios of line PQ joining
P( x1 , y1 , z1 ) and Q( x2 , y2 , z2 ) are x2 x1 a , y2 y1 b and z2 z1 c (say). Then
( x2 x1 ) ( y2 y1 ) ( z2 z1 )
direction cosines are, l ,m ,n , i.e.,
2 2
( x2 x1 ) ( x2 x1 ) ( x2 x1 ) 2
l ( x2 x1 ) PQ , m ( y 2 y1 ) PQ , n ( z2 z1 ) PQ .
If P( x1 , y1 , z1 ) and Q( x2 , y2 , z2 ) are two points in space, then the direction ratios of PQ are
x2 x1 , y2 y1 , z2 z1 . The angle between two lines with direction cosines l1 , m1 , n1 and
l2 , m2 , n2 is given by cos l1l2 m1m2 n1n2 or sin (m1n2 m2 n1 )2
Two lines are perpendicular if l1l2 m1m2 n1n2 0 and parallel if l1 l2 , m1 m2 , n1 n2 or
a1 a2 b1 b2 c1 c2
Projection
Projection of a point on a line: As shown in Fig. 0.19 (a), the projection of a point P on a line
AB is the foot N of the perpendicular PN from P on the line AB . N is also the same point
where the line AB meets the plane through P and perpendicular to AB .
Projection of a segment of a line on another line and its length: As shown in Fig. 0.19 (b), the
projection of the segment AB of a given line on another line CD is the segment AB of CD
where A and B are the projections of the points A and B on the line CD . The length of the
projection AB AN AB cos .
Figure 0.19: (a) Projection of a point on a line; (b) Projection of a segment of a line on another line (c) Projection of
a line joining the two points on another line
Projection of a line joining the two points on another line: As shown in Fig. 0.19 (c), let PQ
be a line segment where P ( x1 , y1 , z1 ) and Q ( x2 , y2 , z2 ) and AB be a given line with d.c.’s
as l , m, n . If the line segment PQ makes angle with the line AB , then projection of PQ is
PQ PQ cos ( x2 x1 ) cos ( y2 y1 ) cos ( z2 z1 ) cos , thus
PQ ( x2 x1 )l ( y2 y1 )m ( z2 z1 )n .
sin 2 1 cos 2 (l12 m12 n12 )(l22 m22 n22 ) (l1l2 m1m2 n1n2 ) 2
sin 2 (l1m2 l2 m1 ) 2 ( m1n2 m2 n1 ) 2 (n1l2 n2 l1 ) 2 .
When d.r.’s of the lines are given if a1 , b1 , c1 and a2 , b2 , c2 are d.r.’s of given two lines, then
( a1b2 a2b1 ) 2
angle between them is given by sin
a12 b12 c12 a22 b22 c22
Condition of perpendicularity: If the given lines are perpendicular, then 900 i.e.
90 o cos 0 l1l2 m1m2 n1n2 0 or a1a2 b1b2 c1c2 0
Condition of parallelism: If the given lines are parallel, then
o 2 2 2
0 sin 0 (l1m2 l2 m1 ) ( m1n2 m2 n1 ) ( n1l2 n2l1 ) 0 , which is true, only
when l1m2 l2 m1 0 , m1n2 m2 n1 0 and n1l2 n2l1 0 l1 l2 m1 m2 n1 n2 .
Similarly, a1 a2 b1 b2 c1 c2 .
Vector form: Let the vector equations of two lines be r a1 b1 and r a2 b 2 . As the lines
are parallel to the vectors b1 and b2 respectively, therefore angle between the lines is same as the
angle between the vectors b1 and b2 . Thus if is the angle between the given lines, then
cos (b1 b 2 ) b1 b2 .
If the lines are perpendicular, then b1 b 2 0 .
If the lines are parallel, then b1 and b2 are parallel, therefore b1 b 2 for some scalar .
x x1 y y1 z z1
Angle between two lines: The angle between the two lines be and
a1 b1 c1
x x2 y y2 z z2
a2
b2
c2
is cos ( a1a2 b1b2 c1c2 ) a12 b12 c12 a22 b22 c22
Condition of perpendicularity: If the lines are perpendicular, then a1a2 b1b2 c1c2 0
Condition of parallelism : If the lines are parallel, then a1 a2 b1 b2 c1 c2 .
Reduction of Cartesian form of the equation of a line to Vector form and vice-versa
x x1 y y1 z z1
Cartesian to Vector: Let the Cartesian equation of a line be
. This is the
a b c
equation of a line passing through the point A( x1 , y1 , z1 ) and having direction ratios a, b, c . In
vector form this means that the line passes through point having position vector
a x1 i y1 j z1 k and is parallel to the vector m a i b j c k . Thus, the vector form of (i) is
r a m or r ( x1 i y1 j z1 k ) (a i b j c k ) , where is a parameter.
Vector to Cartesian: Let the vector equation of a line be r a m , where
a x1 i y1 j z1 k , m a i b j c k and is a parameter. Now put
r x i y j z k , a x1 i y1 j z1 k and m ai b j ck in r am, we get,
x i y j z k ( x1 i y1 j z1 k ) (a i b j c k ) . Equating coefficients of i, j, k , we get
x x1 y y1 z z1
x x1 a , y y1 b , z z1 c or .
a b c
x x1 y y1 z z1
Intersection of two lines: Let the two lines be …(i) and
a1 b1 c1
x x2 y y2 z z2
…(ii). The following steps are for finding the point of intersection:
a2 b2 c2
Step I: Write the co-ordinates of general points on (i) and (ii). The co-ordinates of general points
x x1 y y1 z z1 x x2 y y2 z z2
on (i) and (ii) are given by and ,
a1 b1 c1 a2 b2 c2
respectively, i.e., (a1 x1 , b1 y1 , c1 z1 ) and (a2 x2 , b2 y2 , c2 z2 )
Step II: If the lines (i) and (ii) intersect, then they have a common point. a1 x1 a2 x2 ,
b1 y1 b2 y2 and c1 z1 c2 z2 .
Step III: Solve any two of the equations in and obtained in step II. If the values of and
satisfy the third equation, then the lines (i) and (ii) intersect, otherwise they do not intersect.
Step IV: To obtain the co-ordinates of the point of intersection, substitute the value of (or )
in the co-ordinates of general point (s) obtained in step I.
Length and equation of perpendicular: As shown in Fig. 0.23 (a), the length of the
perpendicular is the distance AP and its equation is the line joining two known points A and P .
The length of the perpendicular is the perpendicular distance of given point from that line.
Reflection or image of a point in a straight line: As shown in Fig. 0.23 (b), if the perpendicular
PL from point P on the given line be produced to Q such that PL QL , then Q is known as
the image or reflection of P in the given line. Also, L is the foot of the perpendicular or the
projection of P on the line.
Vector form
Perpendicular distance of a point
from a line: As shown in Fig. 0.24
(a), let L is the foot of perpendicular
drawn from P ( ) on the line
r a b . Since r denotes the
position vector of any point on the
line r a b . So, let the position Figure 0.24: Perpendicular distance of a point from a line (b) Image
vector of L be a b . Then of a point in a line
(a )b
PL a b (a ) 2 b . The length PL , is the magnitude of PL , and required
|b|
length of perpendicular.
Image of a point in a straight line: As shown in Fig. 0.24 (b), let Q( ) is the image of P in
2(a ) b
r a b . Then, 2a 2 b .
|b|
x x1 y y1 z z1 x x2 y y2 z z2
Cartesian form: Let two skew lines be and
l1 m1 n1 l2 m2 n2
Therefore, the shortest distance between the lines is given by
x2 x1 y2 y1 z2 z1
d l1 m1 n1 ( m1n2 m2 n1 ) 2 ( n1l2 l1n2 ) 2 (l1m2 l2 m1 ) 2 .
l2 m2 n2
Vector form: Let l1 and l2 be two lines whose equations are l1 : r a1 b1 and l2 : r a2 b 2 .
(b1 b 2 ) (a 2 a1 ) [b1 b 2 (a 2 a1 )]
Then, shortest distance PQ
| b1 b 2 | | b1 b 2 |
Shortest distance between two parallel lines: The shortest distance between the parallel lines
| (a 2 a1 ) b |
r a1 b and r a 2 b is given by d .
|b|
Equation of co-ordinate planes: As shown in Fig. 0.26 (a), the equation of XOY plane: z 0 ;
YOZ plane: x 0 ; ZOX plane: y 0
Figure 0.26: (a) Equation of Coordinate plane; (b) Vector equation of a plane; (c) Normal form of a plane
Normal form: As shown in Fig. 0.26 (c), vector equation of a plane normal to unit vector n̂ and
at a distance d from the origin is r nˆ d . If n is not a unit vector, then to reduce the equation
n d d
r n d to normal form we divide both sides by n to obtain r or r nˆ .
|n| |n| |n|
Vector equation of a plane: As shown in Fig. 0.26 (b), vector equation of a plane through the
point A(a) and perpendicular to the vector n is (r a) n 0 or r n a n , which can also be
written as r n d , where d a n . This is known as the scalar product form of a plane.
Equation of a plane passing through a given point and parallel to two given vectors: The
equation of the plane passing through a point having position vector a and parallel to b and c is
r a b c , where and are scalars.
Equation of plane through three points: The equation of plane passing through three non-
x x1 y y1 z z1
collinear points ( x1 , y1 , z1 ) , ( x2 , y2 , z2 ) and ( x3 , y3 , z3 ) is x2 x1 y2 y1 z2 z1 0 .
x3 x1 y3 y1 z3 z1
Equation of plane passing through a given point : Equation of plane passing through the point
( x1 , y1 , z1 ) is A( x x1 ) B( y y1 ) C ( z z1 ) 0 , where A, B, C are d.r.’s of normal to the
plane.
Example 0.33 [AG-2011 (2 marks)]: A plane contains the following three points: P (2,1,5) ,
Q (1, 3, 4) and R(3, 0, 6) . The vector perpendicular to the above plane can be represented as
(a) 2iˆ ˆj kˆ (b) iˆ 2 ˆj 2 kˆ (c) 2iˆ 3 ˆj kˆ (d) iˆ 2 ˆj kˆ
Solution (d): The equation of a plane passing though the given points is given as
x x1 y y1 z z1 x 2 y 1 z 5
x2 x1 y2 y1 z2 z1 0 1 2 3 1 45 0 x 2y z 9, where
x3 x1 y3 y1 z3 z1 32 0 1 65
( x1 , y1 , z1 ) (2,1, 5) , ( x2 , y2 , z2 ) (1, 3, 4) and ( x3 , y3 , z3 ) (3, 0, 6) . As any point on the plane can
be given as r x i y j z k , so the equation of plane can be written as r n d , where
n i 2 j k and d 9 .
Foot of perpendicular from a point to a given plane: Let P be the foot of perpendicular
drawn from the point A( , , ) to a given plane ax by cz d 0 , then AP is parallel to the
x y z
normal to the plane, so its equation is
r (say), on which the point P lies. So
a b c
the point P lies on line and plane thus putting the point ( ar , br , cr ) on the plane we find
the value of r and hence the point P .
Perpendicular distance
Cartesian form : The length of the perpendicular from the point P( x1 , y1 , z1 ) to the plane
ax1 by1 cz1 d
ax by cz d 0 is
a 2 b2 c 2
Vector form : The perpendicular distance of a point having position vector a from the plane
| an d |
r n d is given by p
|n|
The distance between two parallel planes is the algebraic difference of perpendicular
distances on the planes from origin, i.e. distance between two parallel planes
D2 D1
Ax By Cz D1 0 and Ax By Cz D2 0 is
2 2 2
A B C
Example 0.35 [CE-2003 (1 mark)]: If P , Q and R are three points having coordinates (3, 2, 1) ,
(1, 3, 4) , (2,1, 2) in XYZ space, then the distance from point P to plane OQR (O being the origin
of the coordinate system) is given by
(a) 3 (b) 5 (c) 7 (d) 9
Solution (a): The equation of a plane passing though the points is given as
x x1 y y1 z z1 x0 y0 z0
x2 x1 y2 y1 z2 z1 0 1 0 30 40 0 10 x 10 y 5 z 0 , where
x3 x1 y3 y1 z3 z1 20 1 0 2 0
( x1 , y1 , z1 ) (0, 0, 0) , ( x2 , y2 , z2 ) (1, 3, 4) and ( x3 , y3 , z3 ) (2,1, 2) . Now distance from point
10(3) 10( 2) 5( 1) 0
P (3, 2, 1) to the plane 10 x 10 y 5 z 0 is given by 3.
( 10) 2 (10) 2 ( 5) 2
Position of two points w.r.t. a plane: Two points P( x1 , y1 , z1 ) and Q( x2 , y2 , z2 ) lie on the same
or opposite sides of a plane ax by cz d 0 according to ax1 by1 cz1 d and
ax2 by2 cz2 d are of same or opposite signs. The plane divides the line joining the points P and
Q externally or internally according to P and Q are lying on same or opposite sides of the plane.
If a1a2 b1b2 c1c2 0 , then the planes are perpendicular to each other
a1 b1 c1
If , then the planes are parallel to each other
a2 b2 c2
Vector form : An angle between the planes r1 n1 d1 and r2 n 2 d2 is given by
cos (n1 n 2 ) | n1 || n 2 | .
Example 0.36 [AG-2010 (2 marks)]: The angle of intersection between the planes x 3 y 2 z 10
and 2 x 4 y 5 z 0 is
(a) 30o (b) 60o (c) 75o (d) 90o
Solution (d): If is the angle between the planes x 3 y 2 z 10 and 2 x 4 y 5 z 0 , then
(1)(2) ( 3)(4) (2)(5) o
cos 0 90
2 2 2 2 2 2
(1) ( 3) (2) (2) (4) (5)
Example 0.37 [MN-2016 (1 mark)]: Equation of two planes are z 4 and z 4 3 x . The included
angle between the two planes in degrees, is _____.
Solution: The given planes can be written as: 0 x 0 y z 4 …(i), and 3x 0 y z 4 …(ii)
The normal vector of plane (i) can be written as: n1 (0, 0,1) ; The normal vector of plane (i) can be
written as: n2 (3, 0,1) . So the angle between the given two planes is
0 3 0 0 1 1 1 o
cos 1 cos 1 cos 1 ( 0.316) 71.56 o or 108.43 .
2 2 2 2 2 2
0 0 1 ( 3) 0 1 10
Image of a point in a plane: Let P and Q be two points and let be a plane such that (i) Line
PQ is perpendicular to the plane , and (ii) Mid-point of PQ lies on the plane , then either of the
point is the image of the other in the plane , as shown in Fig. 0.27. To find the image of a point in a
given plane, we proceed as follows:
Write the equations of the line passing through P and
x x1 y y1 z z1
normal to the given plane as .
a b c
Write the co-ordinates of image Q as
( x1 ar , y1 br , z1 cr )
Find the co-ordinates of the mid-point R of PQ
Figure 0.27: Image of a point in a plane
Obtain the value of r by putting the co-ordinates of R in
the equation of the plane
Put the value of r in the co-ordinates of Q .
Coplanar lines: Lines are said to be coplanar if they lie in the same plane or a plane can be made to
pass through them.
Every pair of parallel lines is coplanar
Two coplanar lines are either parallel or intersecting
The three sides of a triangle are coplanar
Vector form : If the lines r a1 b1 and r a2 b 2 are coplanar, then [a1b1b 2 ] [a2b1b 2 ]
and the equation of the plane containing them is [r b1 b 2 ] [a1 b1 b 2 ] or [r b1 b 2 ] [a2 b1 b 2 ] .
x x1 y y1 z z1 x x2 y y2 z z2
Cartesian form : If the lines and are coplanar
l1 m1 n1 l2 m2 n2
x2 x1 y2 y1 z2 z1
Then l1 m1 n1 0 . The equation of the plane containing them is
l2 m2 n2
x x1 y y1 z z1 x x2 y y2 z z2
l1 m1 n1 0 or l1 m1 n1 0.
l2 m2 n2 l2 m2 n2
Division by plane : The ratio in which the line segment PQ , joining P( x1 , y1 , z1 ) and
ax by1 cz1 d
Q( x2 , y2 , z2 ) , is divided by plane ax by cz d 0 is 1 .
ax2 by2 cz2 d
Division by co-ordinate planes : The ratio in which the line segment PQ, joining P( x1 , y1 , z1 )
and Q( x2 , y2 , z2 ) is divided by co-ordinate planes are as follows:
By yz plane: x1 x2 By zx plane: y1 y2 By xy plane: z1 z2
Equation of a plane through a given line: If equation of the line is given in symmetrical form as
x x1 y y1 z z1
, then equation of plane is a( x x1 ) b( y y1 ) c( z z1 ) 0 , where a, b, c
l m n
are given by al bm cn 0 .
If equation of lines are given in general form as a1 x b1 y c1 z d1 0 , a2 x b2 y c2 z d2 0
then the equation of plane passing through these line is
(a1 x b1 y c1 z d1 ) (a2 x b2 y c2 z d2 ) 0 ., where is a parameter.
Projection of a line on a plane: If P be the point of intersection of given line and plane and Q
be the foot of the perpendicular from any point on the line to the plane then PQ is called the
projection of given line on the given plane.
x x1 y y1 z z1
Image of line about a plane : Let line is , plane is
a1 b1 c1
a2 x b2 y c2 z d 0 . Find point of intersection (say P ) of line and plane. Find image (say Q )
of point ( x1 , y1 , z1 ) about the plane. Line PQ is the reflected line.
Exercise: 0.5
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. A line makes the same angle with each of the x and z axis. If the angle , which it
2 2
makes with y axis, is such that sin 3sin , then cos 2 equals _____
2. The lines x ay b , z cy d and x ay b , z cy d are perpendicular to each
other, if
(a) aa cc 1 (b) aa cc 1 (c) ac ac 1 (d) ac ac 1
3. x 1 y 1 z 1 x3 yk z
If the line and intersect, then k _____.
2 3 4 1 2 1
4. A line with direction cosines proportional to 2, 1, 2 meets each of the lines x y a z and
x a 2 y 2 z . The co-ordinates of each of the points of intersection are given by
(a) (2a, 3a, 3a) , (2a, a, a) (b) (3a, 2a, 3a) , ( a, a, a)
(c) (3a, 2a, 3a) , ( a, a, 2a) (d) (3a, 3a, 3a ) , ( a, a, a)
5. The length of the perpendicular from the origin to line r (4 i 2 j 4 k ) (3 i 4 j 5 k ) is
_____.
6. The image of point (1, 2,3) in the line r (6 i 7 j 7 k ) (3 i 2 j 2 k ) is
(a) (5, 8,15) (b) (5,8, 15) (c) ( 5, 8, 15) (d) (5,8,15)
7. x2 y 3 z4 x 1 y4 z 5
The line and
are coplanar, if
1 1 k k 2 1
(a) k 0, 1 (b) k 0,1 (c) k 0, 3 (d) k 3, 3
8. The ratio in which the plane x 2 y 3 z 17 divides the line joining the point ( 2, 4, 7) and
(3, 5,8) is
(a) 10 : 3 (b) 3 :1 (c) 3 :10 (d) 10 :1
9. The equation of plane passing through the points (2, 2,1) and (9, 3, 6) and perpendicular to the
plane 2 x 6 y 6 z 1 is
(a) 3 x 4 y 5 z 9 (b) 3 x 4 y 5 z 0 (c) 3 x 4 y 5 z 9 (d) None of these
10. Angle between two planes x 2 y 2 z 3 and 5 x 3 y 4 z 9 is
3 2 19 2 9 2 3 2
(a) cos 1 (b) cos 1 (c) cos 1 (d) cos 1
10 30 20 5
11. Distance between two parallel planes 2 x y 2 z 8 and 4 x 2 y 4 z 5 0 is _____
12. A tetrahedron has vertices at O (0, 0, 0) , A(1, 2,1) , B (2,1,3) and C ( 1,1, 2) . Then the angle
between the faces OAB and ABC will be
19 17
(a) cos 1 (b) cos 1 (c) (d)
35 31 6 2
13. A unit vector perpendicular to plane determined by the points P (1, 1, 2) , Q (2, 0, 1) and
R (0, 2,1) is
2i jk 2i jk 2 i j k 2i jk
(a) (b) (c) (d)
6 6 6 6
14. The reflection of the point (2, 1, 3) in the plane 3 x 2 y z 9 is
26 15 17 26 15 17 15 26 17
26 17 15
(a) , , (b) , , (c) , , , ,
(d)
7 7 7 7 7 7 7 7 7
7 7 7
15. The d.r.’s of normal to the plane through (1, 0, 0) and (0,1, 0) which makes an angle 4 with
plane x y 3 , are
(a) 1, 2,1 (b) 1,1, 2 (c) 1,1, 2 (d) 2,1,1
16. x 1 y 1 zk
Value of k such that the line is perpendicular to normal to the plane
2 3 k
r (2 i 3 j 4 k ) 0 is
(a) 13 4 (b) 17 4 (c) 4 (d) None of these
17. The equation of line of intersection of the planes 4 x 4 y 5 z 12 , 8 x 12 y 13z 32 can
be written as
x y 1 z 2 x y z2 x 1 y 2 z x 1 y 2 z
(a) (b) (c) (d)
2 3 4 2 3 4 2 3 4 2 3 4
18. x3 y 6 z4
The plane which passes through the point (3, 2, 0) and the line is
1 5 4
(a) x y z 1 (b) x y z 5 (c) x 2 y z 1 (d) 2 x y z 5
19. x4 y2 zk
The value of k such that lies in the plane 2 x 4 y z 7 is _____.
1 1 2
20. The distance between the line r (i j 2 k ) (2 i 5 j 3 k ) and the plane
r (2 i j 3 k ) 5 is
5 6 7 8
(a) (b) (c) (d)
14 14 14 14
Answers Keys
Answer Keys: Exercise 0.1
1 2 3 4 5 6 7 8 9 10 11 12 13 14 15
c a c a b a c a b b b c a d d
16 17 18 19 20 21 22 23 24 25 26 27 28 29 30
b 1 a a 1 b b b a d 3 c 7 d d
31 32 33 34 35 36 37 38 39 40 41 42 43 44 45
3 4.5 b d a 12 7 b d d d a b 1 c
46 47 48 49 50 51 52 53 54 55 56 57 58 59 60
7 2 c 1 60 d b d a c 0 9 a b d
61 62 63 64 65 66 67 68 69 70 71 72 73 74 75
0 a b d d b d 4 2 2 3 0 b c a
76 77 78 79 80 81 82 83 84 85 86 87 88 89 90
d c a b a d a b c a d d d b a
91 92 93 94 95 96 97 98 99 100 101 102 103 104 105
b c a d c 2 d b b d a a a d c
106 107 108 109 110 111 112 113 114 115 116 117 118 119 120
2 b c c c 5 a b c 3 c a b a b
121 122 123 124 125 126 127
13 4 a 0 1 1 2
2 0 2 4
[AE – 2017 (2 marks)]: Matrix [ A] 3 2 7 and vector {b} 4 are given. If vector { x} is
3 1 5 5
the solution to the system of equations [ A]{x} {b} , which of the following is true for { x} :
(a) Solution does not exist (b) Infinite solution exist
(c) Unique solution exists (d) Five possible solutions exist
Solution (b): Using Cramer’s rule we have:
2 0 2 4 0 2
3 2 7 2(10 7) 0 2(3 6) 0 ; 1 4 2 7 4(10 7) 0 2(4 10) 0 ;
3 1 5 5 1 5
2 4 2 2 0 4
2 3 4 7 2(20 35) 4(15 21) 2(15 12) 0 ; 3 3 2 4 2(10 4) 0 4(3 6) 0
3 5 5 3 1 5
As 1 2 3 0 . So the given system has infinite solutions.
2 6 x1
[AE – 2017 (2 marks)]: Let matrix [ A] . Then for any non-trivial vector { x} ,
0 2 x2
which of the following is true for the value of K {x}T [ A]{x} :
(a) K is always less than zero (b) K is always greater than zero
(c) K is non-negative (d) K can be anything
T 2 6 x1 x1
Solution (d): K {x} [ A]{x} x1 x2 2 x1 6 x1 2 x2
0 2 x2 x2
[ K ] [2 x12 6 x1 x2 2 x22 ] [2( x1 x2 ) 2 2 x1 x2 ] ; so K can be anything.
[BT – 2017 (2 marks)]: The value of c for which the following system of linear equations
1 2 x c
1 2 y 4 has an infinite number of solutions is _____.
Solution: Using the Cramer’s rule, we must have x y 0 ; so
1 2 c 2 1 c
0 , x 0 0 2c 8 0 c 4 ; y 0 0 4c 0 c 4.
1 2 4 2 1 4
Hence answer is 4.
[CE – 2017 (1 mark)]: The matrix P is the inverse of a matrix Q . If I denotes the identity matrix,
which one of the following options is correct?
(a) PQ I but QP I (b) QP I but PQ I (c) PQ I and QP I (d) PQ QP I
1 1
Solution (c): As P Q ; on pre multiplying with Q , we have QP QQ QP I ; also on post
1
multiplying with Q , we have PQ Q Q PQ I . Thus option (c) is correct.
5 1
[CE – 2017 (2 marks)]: Consider the matrix . Which one of the following statements is
4 1
TRUE for the eigenvalues and eigenvectors of this matrix?
(a) Eigenvalue 3 has a multiplicity of 2, and only one independent eigenvector exists.
(b) Eigenvalue 3 has a multiplicity of 2, and two independent eigenvector exists.
(c) Eigenvalue 3 has a multiplicity of 2, and no independent eigenvector exists.
(d) Eigenvalue are 3 and –3, and two independent eigenvector exists.
5 1 5 1
Solution (a): For eigenvalues of matrix A , we must have A I 0 0
4 1 4 1
(5 )(1 ) 4 0 2 6 9 0 ( 3) 2 0 3, 3 which has a multiplicity of 2.
x1
For eigenvector X of matrix A, we must have AX X , so
x2
5 1 x1 x1 5 x1 x2 3 x1 5 x1 x2 3 x1 2 x1 x2 (i )
4 1 x 3 . From (i) and (ii),
2 x2 4 x1 x2 3x2 4 x1 x2 3x2 2 x1 x2 (ii)
we can say that the given matrix have only one independent eigenvectors. So option (a) is correct.
1 5 3 7 T
[CE – 2017 (2 marks)]: If A and B , AB is equal to
6 2 8 4
38 28 3 40 43 27 38 32
(a) (b) (c) (d)
32 56 42 8 34 50 28 56
3 8 1 5 3 8 1 3 5 7 1 8 5 4 38 28
Solution (a): BT , ABT
7 4 6 2 7 4 6 3 2 7 6 8 2 4 32 56
[Similar question was also asked in MT – 2017 (1 mark)]
n
[CS – 2017 (1 mark)]: Let c1 , …, cn be scalars, not all zero, such that i1 ci ai 0 , where ai are
n
column vectors in R . Consider the set of linear equations Ax b , where A [ a1 ,..., an ] and
n
b i 1 ai . The set of equations has
(a) a unique solution at x J n , where J n denotes a n dimensional vector of all 1
(b) no solution (c) infinitely many solutions (d) finitely many solutions
Solution (c): Since the scalars are not all zero. So the column vectors ai for i 1, 2,3, , n are
n
linearly dependent; so A 0 and b i 1 ai Ax b has infinitely many solutions.
2
[CS – 2017 (2 marks)]: Let u and v be two vectors in R whose Euclidean norms satisfy
|| u || 2 || v || . What is the value of such that w u v bisects the angle between u and v ?
(a) 2 (b) 1 2 (c) 1 (d) 1 2
Solution (a): As w is bisecting the angle between u and v , so cos(angle between w and u) is equal
w u wv
to cos(angle between w and v ) ; i.e. w u 2( w v ) , since || u || 2 || v || ; thus
w u w v
2 2
w u 2( w v) (u v ) u 2{(u v ) v} u ( v u) 2(u v ) 2 v
2 2 2
4 v ( v u) 2(u v) 2 v (4 2 ) v (2 )(u v ) . Now LHS is equal to RHS iff
(2 ) 0 and (4 2 ) 0 , which gives 2 .
[CS – 2017 (2 marks)]: Let A be n n real valued square symmetric matrix of rank 2 with
n n
i1 j 1 Ai2j 50 . Consider the following statements:
(I) One eigenvalue must be in [ 5, 5]
(II) The eigenvalue with the largest magnitude must be strictly greater than 5
Which of the above statements about eigenvalues of A is/are necessarily CORRECT?
(a) Both (I) and (II) (b) (I) only (c) (II) only (d) Neither (I) nor (II)
Solution: Let n 2 , so a 2 2 real valued square symmetric matrix of rank 2 with
n n 5 0
i1 j 1 Ai2j 50 , is given by A 0 5 the eigenvalues are 5, 5, i.e. the largest magnitude is
not greater than 5; also both eigenvalues lies in [ 5, 5] .
n n
Now let n 3 , so a 3 3 real valued square symmetric matrix of rank 2 with i1 j 1 Ai2j 50 , is
5 0 0
given by A 0 5 0 the eigenvalues are 0, 5, 5, i.e. the largest magnitude is not greater than
0 0 0
5; and all the eigenvalues lies in [ 5, 5] .
Hence statement (i) is correct and statement (ii) is not correct.
[CS – 2017 (2 marks)]: If the characteristic polynomial of a 3 3 matrix M over (the set of real
numbers) is 3 4 2 a 30 , a , and one eigenvalue of M is 2, then the largest among the
absolute values of the eigenvalues of M is _____.
Solution: As one eigenvalue of M is 2, so 3 4 2 a 30 0 23 4(22 ) a (2) 30 0
3 2
a 11 . Thus characteristic polynomial is 4 11 30 0 ( 2)( 5)( 3) 0
2,5, 3 . So largest among the absolute values of the eigenvalues of M is 5.
5 10 10 (a) 0
(b) 1
[EC – 2017 (1 mark)]: The rank of the matrix M 1 0 2 is (c) 2
3 6 6 (d) 3
5 10 10 By Applying 1 2 2 By Applying 1 2 2
Solution: M 1 0 2 R1 R1 5 M 1 0 2 R2 R2 R1 M 0 2 0 ,
3 6 6
R3 R3 3 1 2 2 R R R
3 3 1
0 0 0
which has two independent rows, hence its rank is 2.
1 1 0 0 0
0 0 1 1 0
[EC – 2017 (1 mark)]: The rank of the matrix 0 1 1 0 0 is _____.
1 0 0 0 1
0 0 0 1 1
Solution:
1 1 0 0 0 1 1 0 0 0 1 1 0 0 0
0 0 1 1 0 0 1 1 0 0 0 1 1 0 0
By Applying By Applying
0 1 1 0 0 0 0 1 1 0 0 0 1 1 0
R2 R3
R4 R4 R1
1 0 0 0 1 1 0 0 0 1 0 1 0 0 1
0 0 0 1 1 0 0 0 1 1 0 0 0 1 1
1 0 1 0 0 1 0 0 1 0
By Applying
By Applying 0 1 1 0 0 0 1 0 1 0
R1 R1 R3
R1 R1 R2 0 0 1 1 0 0 0 1 1 0
R2 R2 R3
R R R
4 4 2 0 0 1 0 1 R R R 0 0 0 1 1
0 0 0 1 1
4 4 3 0 0 0 1 1
1 0 0 1 0 By Applying 1 0 0 0 0
0 1 0 1 0 R R R 0 1 0 0 0
1 1 4
By Applying
0 0 1 1 0 R2 R2 R4 0 0 1 0 0 , which is in echelon form,
R4 R4
0 0 0 1 1 R3 R3 R4 0 0 0 1 1
0 0 0 1 1 R5 R5 R4 0 0 0 0 0
so rank is total number of non-zero rows, which is 4. So rank is 4.
1 2 3 4 5
5 1 2 3 4
(a) –2.5
[EC – 2017 (1 mark)]: Consider the 5 5 matrix A 4 3 . It is (b)
5 1 2 0
(c) 15
3 4 5 1 2
(d) 25
2 3 4 5 1
given that A has only one real eigenvalue. Then the real eigenvalue of A is
Solution (c): For eigenvalue of matrix A, we have A I 0
1 2 3 4 5
5 1 2 3 4
4 5 1 2 3 0. Applying R1 R1 R2 R3 R4 R5
3 4 5 1 2
2 3 4 5 1
15 15 15 15 15 1 1 1 1 1
5 1 2 3 4 5 1 2 3 4
4 5 1 2 3 0 (15 ) 4 5 1 2 3 0
3 4 5 1 2 3 4 5 1 2
2 3 4 5 1 2 3 4 5 1
15 , which is the real eigenvalue.
3 2 0 1 2
[EE – 2017 (1 mark)]: The matrix A 0 1 0 has three distinct eigenvalues and one of its
1 2 0 3 2
T
eigenvectors is 1 0 1 . Which one of the following can be another eigenvector of A ?
T T T T
(a) 0 0 1 (b) 1 0 0 (c) 1 0 1 (d) 1 1 1
Solution (c): The given matrix is a symmetric matrix and the eigenvectors corresponding to distinct
eigenvalues of a real symmetric matrix are orthogonal, i.e. pairwise dot product is zero. So the vector
T T
given in option (c), i.e. 1 0 1 is orthogonal to the given vector 1 0 1 . Hence option (c) is
correct.
[EE – 2017 (2 marks)]: The eigenvalues of the matrix given below are (a) (0, 1, 3)
0 1 0 (b) (0, 2, 3)
0 0 1 (c) (0, 2,3)
0 3 4 (d) (0,1,3)
Solution (a): For eigenvalue of a matrix A, we must have A I 0 , i.e.
1 0
0 1 0 ( ){( )( 4 ) 3} 0 { 2 4 3} 0 0, 1, 3 .
0 3 4
Alternate Method: We know that sum of the eigenvalues of any square matrix is same as the trace of
that matrix. This property is satisfied by 0, 1, 3 .
[IN – 2017 (1 mark)]: If v is a non-zero vector of dimension 3 1 , then the matrix A v vT has a
rank _____.
Solution: As v is a non-zero vector of dimension 3 1 , so rank( v) 1 . Also rank( vvT ) 1 …(i). As
rank( vvT ) rank( v) 1 …(ii). From (i) and (ii), we can conclude that rank( vvT ) 1 .
[Similar question was also asked in PH – 2017 (2 marks)]
T T
[IN – 2017 (2 marks)]: The angle between two vectors x1 2 6 14 and x 2 12 8 16
in radian is _____.
Solution: The angle between the given vectors is
x x 24 48 14 16 1 o
cos 1 2 0.749 cos (0.749) 41.45
x1 x 2 2 2 2 2 2 2
(2) (6) (14) ( 12) (8) (16)
41.45o ( 180) 0.72 radians.
[MA – 2017 (2 marks)]: Let A [ ai j ] be a 10 10 matrix such that ai j 1 for i j and aii 1 ,
where 0 . Let and be the largest and the smallest eigenvalues of A , respectively. If
24 , then equals _____.
Solution:
1 1 1 1 0 0 1
1 1 1 1 By Applying 0 0 1
A 1 1 1 1 Ci Ci C10 A 0 0 1
where 1 i 9
1 1 1 11010 11010
0 0 1
0 0 1
By Applying
A 0 0 1 , which is an upper triangular
R10 R10 ( R1 R2 R
9 )
0 0 0 10 1010
matrix, whose eigenvalues are the diagonal elements, i.e. , , , (9 times), ( 10) . As 0 , so
lowest eigenvalue is ; and largest eigenvalue is ( 10) . Thus 10 24 7 .
2 0 1 (a) 6
(b) 2
[ME – 2017 (1 mark)]: The product of eigenvalues of the matrix P is P 4 3 3
(c) 6
0 2 1 (d) 2
Solution (b): As the product of eigenvalues of the matrix P is
2 0 1
4 3 3 2(3 6) 0 1(8 0) 2 . So option (b) is correct.
0 2 1
[Similar question was also asked in ME – 2017 (1 mark)]
50 70
[ME – 2017 (2 marks)]: Consider the matrix A whose eigenvectors corresponding to
70 80
70 2 80
eigenvalues 1 and 2 are x1 and x 2 , respectively. The value of x1T x 2 is __.
1 50 70
Solution: If 1 and 2 are eigenvalues of matrix A then 1 2 50 80 130 ; also
T 2 80
x1 x 2 70 (1 50) 70(2 80) (1 50)70 70(1 2 ) 70 80 50 70
70
x1T x 2 70 130 70 80 50 70 0 .
[PI – 2017 (1 mark)]: For two non-zero vectors A and B , if A B is perpendicular to A B , then
(a) the magnitude of A is twice the magnitude of B
(b) the magnitude of A is half the magnitude of B
(c) A and B cannot be orthogonal
(d) the magnitudes of A and B are equal
[TF – 2017 (1 mark)]: If the scalar projection of the vector a 3 i k on the vector
b 2 i 2 j k is 5, then a value of is equal to
(a) 21 (b) 9 (c) 1 (d) 6
a b (3 i k ) (2 i 2 j k ) 6
Solution (b): The scalar projection of a on b is . So
b (2) 2 ( 2) 2 (1) 2 3
6
5 9.
3
Classification of Matrices
Vectors: A matrix having a single column (or row) is called a column (or row) matrix or column
(or row) vector. Its entries are called the components of the vector. It is denoted as
2
a a j ,1 j n . For e.g., is a column vector; and 1 3 513 is a row vector.
3 21
Diagonal Matrix: A square matrix all of whose elements, except those in the leading diagonal, are
zero is called a diagonal matrix, aij 0 , whenever i j . A diagonal matrix of order n n having
3 0
d1 , d 2 , , d n as diagonal elements is denoted by diag d1 , d 2 , , d n . For e.g., A is a
0 5
diagonal matrix of order 2 2 to be denoted by A diag 3 5 .
Scalar Matrix: A diagonal matrix whose all the leading diagonal elements are equal is called a scalar
0, i j
matrix. For a square matrix A aij to be a scalar matrix, aij , where m 0 . For
n n
m, i j
5 0
e.g., A is a scalar matrix of order 2 2 .
0 5
Unit matrix or Identity matrix: A diagonal matrix of order n which has unity for all diagonal
elements, is called a unit matrix of order n and is denoted by I n . For a square matrix A aij to
n n
0, i j 1 0
be a unit matrix, aij . For e.g., I 2 is a unit matrix or order 2 2 .
1, i j 0 1
Triangular matrix: A square matrix in which all the elements below/upper the diagonal are zero is
called upper/lower triangular matrix. Thus, for any square matrix A aij , we have Upper
n n
Triangular matrix if aij 0, i j ; and Lower Triangular matrix if aij 0, i j . [This point was
a d e a 0 0
asked in AG-2011 (1 marks)]. For e.g., 0 b f and d
b 0 are Upper and Lower
0 0 c e
f c
Triangular matrices, respectively. Diagonal matrix is both Upper and Lower Triangular matrix.
Example 1.2 [CS-1994 (1 mark)]: In a compact single dimensional array representation for lower
triangular matrices (i.e. all the elements above the diagonal are zero) of size n n , non-zero elements
(i.e. elements of the lower triangle) of each row are stored one after another, starting from the first
row, the index of the (i, j )th element of the lower triangular matrix in this new representation is:
(a) i j (b) i j 1 (c) j {i (i 1) 2} (d) i { j ( j 1) 2}
Solution (c): Let a n n lower triangular matrix We have to convert A a
ij nn into A ak nn
a
11 0 0
a such that the elements a11 a1 , a21 a2 ,
a 0
is given as, A aij ,
21 22
a22 a3 , a31 a4 , a32 a5 , a33 a6 , and so on.
0
From the given options, option (c) is correct, as
an1 an 2 ann for the option (c) all the criteria are met.
Null matrix: If all the elements of a matrix (square or rectangular) are zero, it is called a null or zero
0 0
matrix. For A aij
m n
to be null matrix, aij 0, i , j . For e.g., , 0 0 are null matrices.
0 0
Matrix Multiplication: Two matrices A and B are conformable for the product AB , if the
number of columns in A (pre-multiplier) is same as the number of rows in B (post-multiplier) [This
point was asked in CE-1997 (1 mark)]. Thus, if A aij and B bij are two matrices of
m n n p
order m n and n p , respectively, then their product AB is of order m p and is defined as,
b1 j
r n b
ain ( i row of A ) ( j
2j th th
AB ij air brj ai1b1 j ai 2b2 j ain brn ai1 ai 2
r 1
bnj
column of B ); where, i 1, 2, , m and j 1, 2, , p
Commutative law does not necessarily hold for matrices. If AB BA , then matrices A and B
are called commutative matrices. If AB BA , then matrices A and B are called anti-
commutative matrices.
Matrix multiplication is associative, i.e., A( BC ) ( AB)C .
Matrix multiplication is distributive w.r.t. addition, i.e., A( B C ) AB AC .
Cancellation law does not necessarily hold, i.e., if AB AC , then in general B C even if
AO.
2
Matrix multiplication A A is represented as A . Thus, An AAA n times.
If AB O , then it is not necessary that at least one of the matrix should be zero. [This point was
0 2 1 0 0 0
asked in AE-2007 (2 marks)]. For e.g. for A , B ; AB while
0 0 0 0 0 0
neither A nor B is null matrix.
If A diag ( a1 , a2 , , an ) and B diag (b1 , b2 , , bn ) , then A B diag ( a1b1 , a2 b2 , , anbn ) .
Thus, A n diag ( a1n , a2n , , a nn ) .
If A and B are diagonal matrices of same order, then AB BA or diagonal matrices (of same
orders) are commutative.
( A B) 2 ( A B )( A B) A2 AB BA B 2 [This point was asked in AE-2007 (2 marks),
ME-2014 (1 mark)]
If A aij and B bij are two matrices of order m n and n p , then total number of
m n n p
multiplication in AB is m n p .
cos sin a 0
Example 1.3 [CS-1996 (2 marks)]: The matrices A and B commute
sin cos 0 b
under multiplication.
(a) If a b or n , n I (b) Always (c) Never (d) If a cos b sin
Solution (a): The two matrices will cos sin a 0 a cos b sin
commute under multiplication if and only AB 0 b a sin b cos
if AB BA b sin a sin sin cos
either a b or sin 0 n , n I a 0 cos sin a cos a sin
BA
. Thus option (a) is correct. 0 b sin cos b sin b cos
Example 1.4 [CS-1997 (2 marks)]: Let a aij be an n rowed square matrix and I12 be the
matrix obtained by interchanging the first and second rows of the n rowed identity matrix. Then
AI12 is such that its first
(a) row is same as its second row (b) row is the same as the second row of A
(c) column is same as the second column of A (d) row is all zero
a11 a12 a1n 1 0 0 0 1 0
a a22 a2 n 0 1 0 1 0 0
Solution (c): Let A and I
21
I12
an1 an 2 ann 0 0 1 0 0 1
a11 a12 a1n 0 1 0 a12 a11 a1n
a a22 a2 n 1 0 0 a22 a21 a2 n
AI12 .
21
an1 an 2 ann 0 0 1 an 2 an1 ann
st nd
Hence 1 column of AI12 is same as 2 column of A .
Example 1.5 [CE-2004 (1 mark)]: Real matrices [ A]31 , [ B ]33 , [C ]35 , [ D]53 , [ E ]55 and [ F ]51 are
given. Matrices [ B ] and [ E ] are symmetric. Following statements are made with respect to these
matrices. (I) Matrix product [ F ]T [C ]T [ B][C ][ F ] is a scalar; (II) Matrix product [ D]T [ F ][ D ] is
always symmetric. With reference these two statements, which of the following applies?
(a) Statement I is true but II is false (b) Statement I is false but II is true
(c) Both the statements are true (d) Both the statements are false
Solution (d): Statement 1 is false, since the matrix product , when it exists, is always another matrix
and not a scalar. Matrix product [ D]T [ F ][ D ] is not possible as number of columns
in [ F ] is not same as number of rows in [ D ] . Hence both statements are false.
Example 1.6 [CE-2005 (1 mark)]: Consider the matrices [ X ]43 , [Y ]43 and [ P ]23 . The order of
T T
P( X T
Y) P
1
will be
(a) 2 2 (b) 3 3 (c) 4 3 (d) 3 4
T T
Solution (a): The order of A P ( X Y ) P T 1
1
((2 3)((3 4)(4 3)) (3 2))T
equal
1 T T T T
A ((2 3)(3 3) (3 2)) ((2 3)(3
3)(3 2)) ((2 3)(3
2)) (2 2) (2 2) .
equal equal
cos sin
Example 1.7 [AE-2007 (1mark)]: If f ( ) , then f ( ) f ( )
sin cos
(a) f ( ) (b) f ( ) (c) f ( ) (d) 2 2 zero matrix
Solution (b):
cos sin cos sin
cos( ) sin( )
f ( ) f ( ) f ( )
sin cos sin
cos sin( ) cos( )
1 1 3 1
Example 1.8 [MN-2007 (2 marks)]: If A B and A B , the value of A B , is
3 0 1 4
1 1 1 1 1 1 1 1 1
(a) 4 (b) 2 (c) (d)
0 3 0 3 0 3 2 0 3
1 1 3 1 4 0 2 0
Solution (b): ( A B ) ( A B) 2A A
3 0 1 4 4 4 2 2
1 1 3 1 2 2 1 1
Also, ( A B ) ( A B) 2B B
3 0 1 4 2 4 1 2
2 0 1 1 2 0 2 0 2 2 1 1
Thus AB 2 .
2 2 1 2 2 2 2 4 0 6 0 3
Example 1.9 [TF-2011 (1 mark)]: X and Y are two matrices such that XY and X Y are both
defined. The CORRECT statement from amongst the following is
(a) X and Y are square matrices of same order (b) X and Y are rectangular matrices
(c) X and Y are diagonal matrices of different (d) X is a square matrix; Y is a rectangular
order matrix
Solution (a): Let matrix X is of order m1 n1 ; and Y is of order m2 n2 . So, for existence of XY
we must have n1 m2 ; also for existence of X Y we must have m1 m2 and n1 n2 . Hence from
the three relations: n1 m2 , m1 m2 and n1 n2 m1 n1 m2 n2 X and Y are square
matrices of same order.
2 1 1 3 0
Example 1.10 [AG-2013 (1 mark)]: If P A B , where A , B 2 1 2 ; then P is
3 0
7 2 4 7 2 2 4 2 4 7
(a) (b) 3 9 0 (c) 3 9 (d) 0 3 9
9 0
Solution (b): As A is a 2 2 and B is a 2 3 matrix, so the product P A B exist whose order is
2 1 1 3 0 2 2 6 1 0 2 4 7 2
2 3 . Hence, P .
3 0 2 1 2 3 0 9 0 0 0 3 9 0
[Similar question was also asked in BT-2013 (1 mark)]
Example 1.11 [CE-2013 (2 marks)]: There are three matrices P (4 2) , Q (2 4) and R(4 1) . The
minimum multiplication required to compute the matrix PQR is …………
Solution: If we first multiply PQ and then ( PQ) R then, the total number of multiplication in PQ is
4 2 4 32 and in ( PQ) R is 4 4 1 16 . So total number of multiplication in PQR is
32 16 48 . If we first multiply QR and then P (QR) then, the total number of multiplication in QR
is 2 4 1 8 and in P (QR) is 4 2 1 8 . So total number of multiplication in PQR is 8 8 16 .
Hence answer is min(48,16) 16 .
2 2 3
e.g., if A 3 2 3 then tr ( A) 2 2 2 6 . [This example was asked in MN-2008 (1 mark)].
4 1 2
If A is an n n identity matrix then aii 1 for all 1 i n , hence tr ( A) 1 1 n - times n .
[This point was asked in MT-2009 (1 mark)].
If A aij , B bij be any two matrices and be any scalar. Then
n n n n
tr ( A) tr ( A) tr ( A B ) tr ( A) tr ( B) tr ( AB ) tr ( BA)
Transpose of Matrix: The matrix obtained from any given matrix A, by interchanging rows and
T
columns, is called the transpose of A and is denoted by A or A . If A aij and AT bij ,
m n n m
1 2
1 4 7
then bij a ji , i , j . For example, if A 4 5 then, AT
2 5 8 23
7 8 32
( AT )T A ( A B )T AT BT
( A)T AT , where, being a scalar AT A
( AB)T BT AT . [This point was asked in CS-1994 (1 mark)].
Proof: Let A aij and B bij be any two matrices. The AB is an m p matrix and
m n n p
T T T T
therefore ( AB)T is a p m matrix. Since A and B are n m and p n , therefore, B A is a
T T T
p m matrix. Thus, the two matrices ( AB) and B A are of the same order such that,
r n r n r n
( AB)
T
ij
( AB) ji a jr bri bri a jr B A B
T
ir
T
rj
T
A
T
ij
. Hence, by the definition
r 1 r 1 r 1
T T
of equal matrices, ( AB) BT AT = B A . In general, ( ABC )T C T B T AT .
T
Example 1.13 [CS-1994 (1 mark)]: Let A and B be real symmetric matrices of size n n . Then
which one of the following is true?
T
(a) AA 1 (b) A A
1
(c) AB BA (d) ( AB )T BA
Solution (d): As, A and B be real symmetric matrices A A
T
and BT B .
( AB )T BT AT ( AB )T BA
3 2 1 1
Example 1.15 [CE-2014 (1 mark)]: Given the matrices J 2 4 2 and K 2 , the product
1 2 6 1
T
K JK is _____.
3 2 1 1 1
Solution: K JK 1 2 1 2 4 2 2 8 12 11 2 43 . Hence answer is 43.
T
1 2 6 1 1
[Similar question was also asked in CE-2001 (2 marks)]
3 1 x
Example 1.16: [EE-2016 (2 marks)]: Let P . Consider the set S of all vectors such
1 3 y
a x
that a 2 b 2 1 , where P . Then S is
b y
(a) a circle of radius 10 (b) a circle of radius 1 10
T T
(c) an ellipse with major axis along 1 1 (d) an ellipse with minor axis along 1 1
a x a 3 1 x a 3 x y a 3 x y
Solution (d): P
b y b 1 3 y b x 3 y b x 3 y
a 2 b 2 1 (3 x y ) 2 ( x 3 y ) 2 1 10 x 2 12 xy 10 y 2 1 0 …(i), which is 2nd degree curve.
On comparing (i) with ax 2 2hxy by 2 2 gx 2 fy c 0 , we have a 10 , h 6 , b 10 ,
a h g 10 6 0
g f 0, c 1 . As h b f 6 10 0 1(100 36) 64 0 ; also
g f c 0 0 1
2 2 2
h ab 6 (10)(10) 64 h ab 0 . So equation (i) is an ellipse. For finding the centre of
the ellipse (i), we have to find the point of intersection of the equations which are found by partial
derivatives of (i) w.r.t. x and y . Let f ( x, y ) 10 x 2 12 xy 10 y 2 1 ; so f x 20 x 12 y 0 , and
f y 12 x 20 y 0 ; on solving these equations we get x 0 , y 0 . Thus the major and minor
axis of the ellipse (i) passes through the centre of ellipse (i), i.e. (0, 0) .
Because of the term ‘ xy ’ in (i), the major and minor axis of the ellipse are not parallel to x or y
axis. To remove the ‘ xy ’ term from (i), we need to rotate the x y axis by an angle
1 2h 1
2(6) 1
tan 1 tan 1
. Hence we can conclude that the major or minor axis of
2 ab 2 10 10 2 2 4
o o
the ellipse (i) is at an angle of 45 with the ve x axis; or in other words if major axis is at 45 with
the ve x axis then minor axis is at (90 45)o 135o with the ve x axis or vice-versa. So
o
equation of line passing through (0, 0) making an angle of 45 with the ve x axis is y x 0
o
…(ii); and equation of line passing through (0, 0) making an angle of 135 with the ve x axis is
y x 0 …(iii).
The point of intersection of (ii) with (i) is A and B , which is given by putting (ii) in (i)
10 x 2 12 x 2 10 x 2 1 0 x 1 4 2 y 1 4 2 .
2 2
Thus A(1 4 2 ,1 4 2) and B ( 1 4 2 , 1 4 2) ; and so AB (2 4 2) (2 4 2) 1 2 .
Similarly, point of intersection of (iii) with (i) is C and D , which is given by putting (iii) in (i)
10 x 2 12 x 2 10 x 2 1 0 x 1 2 2 y 1 2 2 . Thus C (1 2 2 , 1 2 2) and
2 2
D ( 1 2 2 ,1 2 2) ; and so CD (2 2 2) (2 2 2 ) 1 .
As AB CD , so the line CD is the major axis whose equation is given (iii), i.e.
y x 0 y x ; and the line AB is the minor axis whose equation is given (ii), i.e.
y x 0 y x . Hence option (d) is correct.
Conjugate of Matrix: The matrix obtained from any given matrix A containing complex number
as its elements, on replacing its elements by the corresponding conjugate numbers is called conjugate
1 2i 2 3i 1 2i 2 3i
of A and is denoted by A . For example, if A then, A .
4 5i 5 6i 4 5i 5 6i
( A) A If A A A is purely imaginary matrix
1 2i 2 3i 1 2i 4 5i
if, A , then, A .
4 5i 5 6i 2 3i
5 6i
( A ) A ( A B ) A B ( kA) kA , k or ( AB) B A
0 2 3
For e.g., A 2 0 4 is a skew-symmetric matrix as A A . [This example was asked in
T
3 4 0
AG-2012 (1 mark)]
T n 1 T
If A is a symmetric matrix, then, A, kA, A , A , A , B AB are also symmetric matrices, where,
n N , k R and B is a square matrix of order that of A .
If A is a skew-symmetric matrix, then
2n
A is a symmetric matrix for n N kA is a skew-symmetric matrix for k R
A2 n 1 is a skew-symmetric matrix for BT AB is skew-symmetric matrix, where B is
nN a square matrix of order that of A
a b ic 0 2 i
if aij aij , i, j A A . For example, A and A are
b ic d 2 i 0
Hermitian and skew-Hermitian matrix, respectively.
If A is a Hermitian matrix, then aii aii aii is real i . Thus every diagonal element of a
Hermitian matrix must be real. A Hermitian matrix over a set of real numbers is actually a real
symmetric matrix.
If A is a skew-Hermitian matrix, then, aii aii aii aii 0 , i.e., aii must be purely
imaginary or zero.
A skew-Hermitian matrix over the set of real numbers is actually a real skew-symmetric matrix.
l 0 sin
Example 1.17 [XE-2008 (2 marks)]: The matrix A 0 1 m is orthogonal, if
n 0 cos
(a) l sin , m cos , n 0 (b) l sin , m 0, n cos
(c) l cos , m sin , n 0 (d) l cos , m 0, n sin
l 0 sin l 0 n 1 0 0
Solution (d): For orthogonal matrix, AA I 0 1 m 0 0 1 0
T
1 0
n 0 cos sin m cos 0 0 1
Nilpotent Matrix: A square matrix A is called a nilpotent matrix if there exist a positive integer m
m1 m 2
such that Am O and A O, A O , A O . If m is the least positive integer such that
m
A O , then m is called the index of the nilpotent matrix A .
Permutation matrices: It is a square matrix whose entries are all 0’s and 1’s, with exactly single
0 1 0
0 1
‘1’ in each row and exactly single ‘1’ in each column. For e.g. 0 0 1 are permutation
1 0
1 0 0
matrices. Equivalently, permutation matrix is the identity matrix I with its rows rearranged. A
T
permutation matrix is non-singular, and satisfies AA I , thus the determinant is always 1 .
Example 1.18 [ME-2004 (2 marks)]: If matrix Solution (a): If given matrix to be singular,
8 x 0 8 x 0
4 0 2 is singular, find value of x . 4 0 2 0 8(0 12) x (0 24) 0
12 6 0 12 6 0
(a) 4 (b) 6 (c) 8 (d) 12 x4
Example 1.19 [MT-2013 (2 marks)]: Which one of the following attributes is NOT correct for the
cos sin 0
matrix? sin cos 0 , where 60o
0 0 1
Minors and Cofactors: If A is a square matrix, then the minor of the entry in the ith row and j th
column is the determinant of the sub-matrix formed by deleting the ith row and j th column. This
number is often denoted by M ij . The corresponding cofactor ( Cij ) is obtained by multiplying the
i j
minor by ( 1)i j , i.e., Cij (1) M ij .
a d g
e h b h b e
For, e.g., if A b e h , then M 11 ; M 12 ; M 13 ; and so on.
f i c i c f
c f i
Also, C11 ( 1) M 11 ; C12 ( 1)1 2 M 12 ; C13 (1)13 M 13 and so on.
11
The value of the determinant is the sum of product of the elements of any row or column with the
a11 a12 a13
corresponding cofactors. For e.g., A a21 a22 a23 then, in terms of notation of the cofactors,
a31 a32 a33
A a11C11 a12 C12 a13C13 a21C21 a22C22 a23C13 a31C31 a32C32 a33C33
A a11C11 a21C 21 a31C31 a12 C12 a22C 22 a32C32 a13C13 a23C23 a33C33
The sum of product of the elements of any row or column with the cofactors of the corresponding
elements of any other row or column is zero, i.e.,
a11C21 a12C22 a13C23 0; a11C31 a12 C32 a13C33 0
If A1 , A2 , , An are square matrix of same order, then A1 A2 An A1 A2 An
If k is scalar, then kA k n A , where, n is the order of matrix A . [This point was asked in
CE-1999 (1 mark), MT-2012 (1 mark)]
If A and B are square matrices of same order then AB BA even if AB BA
The determinant of a triangular matrix is the product of all elements along principle diagonal.
If A and B are two n n square matrices then det( AB ) (det A)(det B) . [This point was asked
in EC-2014 (1 mark)]
If A and B are two n n square matrices and AB O (i.e., AB is a null matrix) then
det( AB ) 0 (det A)(det B) 0 and so we have the following cases: (i) A 0 and B 0 (ii)
A 0 and B 0 (iii) A 0 and B 0 . [This point was asked in IN-2010 (2 marks)].
Example 1.23 [AE-2016 (1 mark)]: If A and B are both non-singular n n matrices, then which of
the following statement is NOT TRUE. Note: det represent the determinant of a matrix.
(a) det( AB) det( A ) det(B) (b) det( A B) det( A ) det(B)
(c) det( AA 1 ) 1 (d) det( AT ) det( A)
Solution (b): For any two nn matrices A and B , det( AB) det( A ) det(B) ;
1 1
det( AA ) det( A) det( A ) det( A) det( A ) 1 ; det( AT ) det( A) . But
det( A B) det( A ) det(B) .
Example 1.24 [BT-2016 (2 marks)]: The Solution: Since the given matrix is triangular (upper)
value of determinant A given below matrix, hence its determinant is the product of all
5 16 81 elements along principle diagonal, which is equal to
A 0 2 2 is _____. A 5(2 16 2 0) 160 .
[Similar question was also asked in CS-1997 (1 mark)]
0 0 16
a 0 3 7
2 5 1 3
Example 1.25 [EC-2016 (2 marks)]: The matrix A has det( A) 100 and trace
0 0 2 4
0 0 0 b
( A) 14 . The value of a b is ___.
Solution: Expanding the given determinant along first column,
5 1 3 0 3 7
det( A) a 0 2 4 2 0 2 4 10ab , as det( A) 100 10ab 100 ab 10 …(i).
0 0 b 0 0 b
sum of diagonal elements of a square matrix gives trace of that matrix, so
a 5 2 b 14 a b 7 …(ii).
From (i) and (ii), we have ( a b ) 2 ( a b ) 2 4 ab 7 2 4(10) 49 40 9 a b 3 .
j 1 i 1
an1 an 2 ann an1 an 2 ann
For easier calculations, we shall expand the determinant along that row or column which contains
maximum number of zeros.
Sarrus rule for Expansion (Valid only for determinant of order 3):
The three diagonals sloping down to the right give the three positive terms
and the three diagonals sloping down to left the three negative terms.
a1 b1 c1
a2 b2 c2 a1b2 c3 b1c2 a3 c1a2 b3 a3b2 c1 b3 c2 a1 c3 a2 b1
a3 b3 c3
th
For n order determinant we have n! number of terms in its expansion. [This point was asked
in CE-1999 (1 mark)]
0 0 0 0 0 1
0 0 0 0 1 0
0 0 0 1 0 0
Example 1.26 [CS-2014 (2 marks)]: Consider the matrix J 6 which is
0 0 1 0 0 0
0 1 0 0 0 0
1 0 0 0 0 0
obtained by reversing the order of the columns of the identity matrix I 6 . Let P I 6 J 6 , where
is a non-negative real number. The value of for which det( P ) 0 is ……………
Solution:
1 0 0 0 0 0 0 0 0 0 0 1 0 0 0 0
0 1 0 0 0 0 0 0 0 1 0
0 0 0 0 0
0 0 1 0 0 0 0 0 0
0 0 0 0 1 0 0
P I6 J6
0 0 0 1 0 0 0 0 0
0 0 0 0 1 0 0
0 0 0 0 1 0 0 0 0
0 0 0 0 0 1 0
0 0 0 0 0 1 0 0 0
0 0 0 0 0 0 1
Now the det( P ) 0 if and only if any two row or column are same which results 1.
1 0 1 0
4 7 0 2
Example 1.27 [BT-2011 (2 marks)]: Value of the determinant is
1 1 1 1
2 0 2 1
(a) 24 (b) –30 (c) –24 (d) –10
Solution (c): Expanding along first row we get,
7 0 2 4 7 2
1 1 1 1 0 1 1 1 1 7( 1 2) 0 2(2 0) 2(7 2) 0 1(4 7) 24
0 2 1 2 0 1
[Similar questions were also asked in CS-2000, CE-2014 (1 mark)]
f1 ( r ) f 2 (r ) f3 (r ) f1 (r ) f 2 (r ) f3 (r )
r n r 1 r 1 r 1
7.
r 1
r a b c , if r a b c , where, f1 ( r ), f 2 ( r ), f3 ( r )
d e f d e f
q q q
p f1 ( x) p f 2 ( x) p f2 ( x) f1 ( x ) f 2 ( x) f3 ( x)
q
Also p ( x ) a b c , if ( x ) a b c where,
d e f d e f
Differentiation of a Determinant:
Let ( x ) be a determinant of order two. If we write ( x ) C1 C2 , where, C1 and C 2 denotes
the first and second column then, ( x ) C1 C2 C1 C 2 , where, Ci denotes the column
which contains the derivative of all the functions in the ith column.
R1 R1 R1
In similar fashion, if we write, ( x ) , then ( x ) , where, Ri denotes
R2 R2 R2
the row which contains the derivative of all the functions in the ith row.
Let ( x ) be a determinant of order three. If we write ( x ) C1 C2 C3 , where, C1 , C 2 and
C3 denotes the first, second and third column then,
( x) C1 C2 C3 C1 C2 C3 C1 C2 C3 , where, Ci denotes the column which
contains the derivative of all the functions in the ith column.
R1 R1 R1 R1
In similar fashion, if we write, ( x ) R2 , then ( x) R2 R2 R2 , where, Ri
R3 R3 R3 R3
denotes the row which contains the derivative of all the functions in the ith row.
Example 1.36 [TF-2012 (2 marks)]: Consider the following Assertion [A] and Reason [B].
1 0 0
A. M is an orthogonal matrix, but not a skew-symmetric matrix. M 0 cos sin
0 sin cos
T 1 T
B. Because M M and M M
Determine the correctness or otherwise of the above Assertion [A] and Reason [B]
(a) A is right B is wrong (b) A is right B is right
(c) A is wrong B is right (d) A is wrong B is wrong
T
Solution (b): If M is an orthogonal matrix then MM I
1 0 0 1 0 0 1 0 0
MM T 0 cos sin 0
cos sin 0 1 0 I MM T I M T M 1 .
0 sin cos 0 sin cos 0 0 1
1 0 0
M T 0 cos sin M M is not skew-symmetric matrix. So (A), (B) both are correct.
0 sin cos
Inverse of a square Matrix: A non-singular square matrix of order n is invertible if there exist a
square matrix B of same order such that AB I n BA ; in such a case we say that the inverse of A
1
is B and write A B [This point was asked in MN-2016 (1 mark)]. From the properties of
1
adjoint and inverse of a matrix A1 adj ( A) A . For A to be defined, A 0 A must be non-
singular matrix.
1 0 1
Example 1.38 [CS-1994 (2 marks)]: Find the inverse of the matrix A 1 1 1 .
0 1 0
Solution: Since A 1(0 1) 0 1( 1 0) 2 0 , hence A is invertible.
Let Cij be cofactor of aij in A . Then the cofactors of elements of A are given by:
1 1 1 1 1 1 0 1 1 1
C11 1 , C12 0 , C13 1 , C21 1 , C22 0,
1 0 0 0 0 1 1 0 0 0
1 0 0 1 1 1 1 0
C23 1 , C31 1 , C32 2 , C33 1.
0 1 1 1 1 1 1 1
T
1 0 1 1 1 1 1 2 1 2 1 2
adjA
adjA 1 0 1 0 1
0 2 A 0 0 1
A
1 2 1 1 1 1 1 2 1 2 1 2
[Similar questions were also asked in CE-1997 (1 mark), MN-2012 (2 marks)]
T
1 1 1 1 0 T
Now we have to find P as, P (1 0) 1 and adjA Cij
0 1 1 1
1 adjA 1 0 1 0 b11 b12 b11 b12
P . D P 1 B
A 1 1 1 1 b21 b22 b11 b21 b12 b22
Reversal Law for adjoint: If A and B are non-singular square matrices of the same order, then
adj ( AB ) ( adjB )(adjA) ( AB ) 1 B 1 A1
If A is an invertible square matrix, then adj ( AT ) ( adjA)T ( AT ) 1 ( A1 )T
n 2
If A is a non-singular square matrix, then adj ( adjA) A A
Proof: B ( adjB ) B I n for every square matrix of order n . Replacing B by adjA , we get,
( adjA)[ adj ( adjA)] adjA I n A
n 1
I n A(adjA)[adj ( adjA)] A A n 1
In
n 1
( AadjA)[ adj (adjA)] A ( AI n )
n 1
A I n [adj ( adjA)] A A AI n A & AadjA A I n
n 1 n2
A [ I n adj ( adjA)] A A [ adj ( adjA)] A A
1 1
If A is a non-singular matrix, then A A A1 1 A
1 1
Proof: A 0 A1 exist such that AA I A A AA1 I 1
A A1 1 A1 1 A A 0
Inverse of k th power of A is the k th power of the inverse of A , i.e., ( A1 ) k ( Ak ) 1
( Ak ) 1 ( A A A) 1 A1 A1 A1 ( A1 ) k .
2 0.1 1 2 a
Example 1.43 [EC-2005 (2 marks)]: Let A and A1 . Then ( a b)
0 3 0 b
(a) 7/20 (b) 3/20 (c) 19/60 (d) 11/20
Solution (a):
1 2 0.1 1 / 2 a 1 0 1 2a 0.1b 1 0 2a 0.1b 0 (i)
AA I
0 3 0 b 0 1
0 3b 0 1 3b 1 (ii)
(ii) b 1 3 and thus (i) a 0.1 6 a b 1 3 0.1 6 7 20 .
1
Example 1.47 [CE-2008 (1 mark)]: The product of matrices PQ P is
(a) P
1
(b) Q 1 (c) P 1Q 1 (d) PQP 1
Solution (b): [ PQ]1 Q 1 P 1 [ PQ]1 P Q 1 P 1 P Q 1 I Q 1 .
3 5 4 5
Example 1.48 [ME-2009 (1 mark)]: For a matrix M , the transpose of the matrix is
x 3 5
T 1
equal to the inverse of the matrix M M . The value of x is given by
(a) –4/5 (b) –3/5 (c) 3/5 (d) 4/5
3 5 4 5 3 5 x 1 0
Solution (a): M T M 1 MM T I
x 3 5 4 5 3 5 0 1
1 3 x 5 12 25 1 0 3 12 4
0 1 5 x 25 0 x 5 .
3x 5 12 25 1
Example 1.49 [TF-2009 (2 marks)]: If the determinant and trace of a 2 2 matrix M are –1 and 0
1
respectively, then the determinant of M M is
(a) –4 (b) –2 (c) –1 (d) 0
a b
Solution (a): Let the matrix M . As tr ( A) 0 a d 0 d a
c d
a b adj ( M )
M . Now, M 1 adj ( M ) as det( M ) 1 .
c a M
C11 ( 1)11 M 11 a , C12 (1)1 2 M 11 c , C21 ( 1) 2 1 M 21 b , C 22 (1) 2 2 M 21 a
T
T a c a b a b 2a 2b
adj ( M ) Cij M 1 M M 1
b a c a c a 2c 2a
M M 1 2 M det( M M 1 ) det(2M ) 22 det( M ) 4 1 4 .
3 2i i 2
Solution (b): A (3 2i )(3 2i) i 9 4 1 12 0 A is invertible. Now,
i 3 2i
11
C11 ( 1) M 11 3 2i , C12 ( 1)1 2 M 12 i , C 21 ( 1) 21 M 12 i , C22 (1) 2 2 M 12 3 2i
T
T3 2i i 3 2i i adj ( A) 1 3 2i i
adj ( A) Cij A 1
i 3 2i i 3 2i A 12 i 3 2i
4
Example 1.51 [EC-2016 (1 mark)]: Let M I , (where I denotes the identity matrix) and M I ,
M 2 I and M 3 I . Then, for any natural number k , M 1 equals:
4 k 1 4k 2 4 k 3 4k
(a) M (b) M (c) M (d) M
4 1 4 1 1
Solution (c): Multiplying both sides of given relation M I with M , we get M M M
1 3 1 3 4k 3 k
M 3 ( M ) M 1 M 1 M 3 M 1 M 1 M 3 M M M IM I M ( I I ),
where k is any natural number. So M 1 I k M 3 ( M 4 ) k M 3 M 4 k 3 , thus option (c) is correct.
1 1 2 1
Example 1.52 [TF-2016 (1 mark)]: Let A . The determinant of A is equal to
1 2 1
(a) 1 2 (b) 4 3 (c) 3 4 (d) 2
Solution (b): A 1 (1 2)(1 2) 3 4 . So A1 1 A 4 3 .
Algorithm for Finding the Inverse of a Non Singular Matrix by Elementary Row
A In A
Transformations: Let A be non-singular matrix of order n . Write , then perform a
I
sequence of elementary row operations successively on A on the LHS and the pre factor n on
I BA A1 B
the RHS till we obtain the result n , . The following steps will be helpful to find
the inverse of a square matrix of order 3 by using elementary row transformations.
Step I: Introduce unity at the intersection of first row and first column either by interchanging
two rows or by adding a constant multiple of elements of some other row to first row.
Step II: After introducing unity at (1,1) place introduce zeros at all other places in first column.
Step III: Introduce unity at the intersection of 2nd row and 2nd column with the help of 2nd and 3rd
row.
Step IV: Introduce zeros at all other places in the second column except at the intersection of 2nd
row and 2nd column.
Step V: Introduce unity at the intersection of 3rd row and third column.
Step VI: Finally introduce zeros at all other places in the third column except at the intersection
of third row and third column.
5 0 2
Example 1.53 [EE-1998 (2 marks)]: If A 0 3 0 . The inverse of A is
2 0 1
1 0 2 5 0 2 1 5 0 1 2 15 0 1 2
(a) 0 1 3 0 (b) 0 1 3 0 (c) 0 1 3 0 (d) 0 13 0
2 0 5 2 0 1 1 2 0 1 1 2 0 1
5 0 2 1 0 0 1 0 2 5 1 5 0 0
By Applying
Solution (a): A IA 0 3 0 0 1 0 A 0 3 0 0 1 0 A
R1 R1 5
2 0 1 0 0 1 2 0 1 0 0 1
1 0 2 5 1 5 0 0 By Applying 1 0 2 5 1 5 0 0
By Applying
0 3 0 0 1 0 A R2 R2 3 0 1 0 0 1 3 0 A
R R 2 R
3 3 1 0 0 1 5 2 5 0 1 R R (1 5) 0 0 1 2 0 5
3 3
1 0 0 1 0 2 1 0 2
By Applying 0 1 0 0
1 3 0 A. I A A A 0 1 3 0 .
1 1
R R (2 5) R
1 1 3 0 0 1 2 5
0 2 0 5
Similar questions were also asked in EE-1995, AG-2009, CH-2010 (1 mark), MN-2008, CE-2007,
PI-2008 (2 marks)
Exercise 1.1
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
0 1 0 2
1. If A and B , then value of for which A B is (a) 1 (b) 2
1 1 5 1 (c) 2 (d)
2 3
2. If A and A 125 then maximum value of is _____.
2
2 x 1 3y x 3 y 2 2
3. If A and B , then for what values of x and y s.t. A B ?
y 5 y
2
0 0 6
1 1 1 4 2 2
28. Let A 2 1 3 and 10 B 5 0 . If B is the inverse of A , then _____.
1 1 1 1 2 3
2 1
29. If there are three square matrices A, B, C of same order satisfying the equation A A and let
n ( n2)
B A2 and C A2 , where n N , then B C _____.
30. The inverse of the matrix 1 2 1 2 1 2 1 2 1 2 1 2
0 1 2 (a) 4 3 1 (b) 4 3 1
A 1 2 3 is
5 2 3 2 1 2 5 2 3 2 1 2
3 1 1 1 2 1 2 1 2 1 2 1 2 1 2
(c) 4 3 1 (d) 4 3 1
5 2 3 2 1 2 5 2 3 2 1 2
31. If A aij , B bij and C cij , then which one of the following is not defined?
33 32 31
9
43. If A is a square matrix of order n such that det{adj ( adjA)} {det( A)} , then n _____.
44. If A , B are symmetric matrices of same order and X AB BA , Y AB BA , then ( XY )T
(a) XY (b) YX (c) YX (d) None of these
45. If A is a 3 3 skew-symmetric matrix, then tr ( A)
(a) 1 (b) –1 (c) det( A) (d) None of these
Consistent system: If the system of the equation (1.2) has one or more solutions, then it is said
to be a consistent system of equation.
If it has unique solution then consistent system is determinate.
If it has more than one solution then consistent system is indeterminate.
Inconsistent system: If the system of the equation (1.2) has no solutions, then it is said to be an
inconsistent system of equation.
Over-determined system: If the system of equation (1.2) has more equations than unknowns,
i.e., m n , then it is said to be an over-determined system. In general, an over-determined system
may be inconsistent (i.e., having no solution) or consistent (having unique or infinite number of
solutions). [This point was asked in CE-2005 (1 mark)]
Determined system: If the system of equation (1.2) has equal number equations and unknowns,
i.e. m n , then it is said to be an determined system. In general, a determined system may be
inconsistent (i.e., having no solution) or consistent (having unique or infinite number of
solutions).
Under-determined system: If the system of equation (1.2) has less equations than unknowns,
i.e., m n , then it is said to be an under-determined system. In general, an under-determined
system may be inconsistent (i.e., having no solution) or consistent (having infinite number of
solutions).
Cramer’s Rule: Cramer’s rule is applicable when number of equations is same as the number of
unknowns, i.e., for (1.5), we have m n . Cramer’s rule express the solution in terms of the
determinants of the coefficient matrix and of matrices obtained from it by replacing one column by
the vector of right hand sides of the equations. Consider the system of equation given by (1.5), then
a11 a12 a1n b1 a12 a1n a11 a12 b1
a21 a22 a2 n b2 a22 a2 n a21 a22 b2
Let, and 1 ; , n , where i
an1 an 2 ann bn an 2 ann an1 an 2 bn
is obtained by replacing the elements of i th column by the elements of B . Then, the solution of
system of non-homogeneous linear equations is given by,
x1 1 ; x2 2 ; xn n (1.7)
Example 1.55 [CE-2007 (2 marks)]: For what values of and the following simultaneous
equations have an infinite number of solutions? x y z 5 ; x 3 y 3 z 9 ; x 2 y z .
(a) 2, 7 (b) 3, 8 (c) 8, 3 (d) 7, 2
Solution (a): For infinitely many solutions, 1 2 3 0
1 1 1 5 1 1
We have, 1 3 3 2 4 0 2; 1 9 3 3 6 12 0 2;
1 2 2
1 5 1 1 1 5
2 1 9 3 4 2 6 0 7; 3 1 3 9 2 14 0 7
1 1 2
2
Example 1.56 [IN-2007 (1 mark)]: Let A be n n real matrix such that A I and y be an n
dimensional vector. Then the linear system of equations Ax y has
(a) no solution (b) more than one but finitely many independent solutions
(c) a unique solution (d) infinitely many independent solutions
Solution (c): A2 I A2 I AA 1 A A 1 A 1 0 . Hence 0 . Also for any
value of 1 , 2 and 3 , we can say that the system Ax y has a unique solution.
Example 1.59 [ME-2008 (2 marks)]: For what values of a , if any, will the following system of
equations in x, y and z have a solution? 2 x 3 y 4 ; x y x 4 ; x 2 y z a
(a) any real number (b) 0 (c) 1 (d) There is no such value
Solution (a): As the given system of equations are: 2 x 3 y 4 ; 2 x y 4 ; x 2 y z a
2 3 0
For finding the nature of solution, first we have to find , i.e., 2 1 0 1(2 6) 4
1 2 1
0 . So, it has a unique solution for any value of a .
Example 1.61 [BT-2013 (2 marks)]: The solution of the following set of equations: x 2 y 3 z 20
; 7 x 3 y z 13 ; x 6 y 2 z 0 ; is
(a) x 2, y 2, z 8 (b) x 2, y 3, z 8 (c) x 2, y 3, z 8 (d) x 8, y 2, z 3
Solution (b): We have,
1 2 3 20 2 3 1 20 3 1 2 20
7 3 1 91; 1 13 3 1 182; 2 7 13 1 273; 3 7 3 13 728
1 6 2 0 6 2 1 0 2 1 6 0
1 182 2 273 3 728
x 2; y 8. 3; z
91 91 91
[Similar questions were also asked in PI-2009, MN-2016 (2 marks), AG-2014 (1 mark)]
Example 1.62 [AG-2014 (2 marks)]: Consider the following set of linear equations x1 x2 x3 6 ,
2 x1 2 x2 3 x3 14 , 3 x1 x2 2 x3 14 . The solution for this set exists only when the value of x2 is
_____.
1 1 1 6 1 1 1 6 1 1 1 6
Solution: 2 2 3 2; 1 14 2 3 6; 2 2 14 3 2; 3 2 2 14 4
3 1 2 14 1 2 3 14 2 3 1 14
1 6 2 2 3 4
x1 3; x2 1; x3 2 . Hence answer is 1.
2 2 2
1 2 0 1
Example 1.63 [XE-2014 (2 marks)]: Let r , s . If A 2 0 3 and b 1 , then the
r s 0 s 1
system of linear equations AX b has
(a) no solutions for s 2r (b) infinitely many solutions for s 2r 2
(c) a unique solution for s 2r 2 (d) infinitely many solutions for s 2r 2
1 2 0
Solution (d): For a given linear system of equations: 2 0 3 3(2r s ) ,
r s 0
1 2 0 1 1 0 1 2 1
1 1 0 3 3( s 2) , 2 2 1 3 3( r s 1) , 3 2 0 1 3s 2r 4
s 1 s 0 r s 1 0 r s s 1
Now For a unique solution: 0 2r s 0 2r s . For no solution:
0 2r s 0 2r s and any one of 1 or 2 or 3 is non-zero. For infinitely many
solutions 1 2 3 0 s 2r and s 2 .
2 1 3 a 5
Example 1.64 [EC-2014 (2 marks)]: The system of linear equations 3 0 1 b 4 has
1 2 5 c 14
(a) a unique solution (b) infinitely many solutions
(c) no solution (d) exactly two solutions
Solution (b): As we have non-homogeneous equation with
2 1 3 5 1 3
3 0 1 4 14 18 0 ; 1 4 0 1 10 34 24 0 ;
1 2 5 14 2 5
2 5 3 2 1 5
2 3 4 1 68 70 138 0 ; 3 3 0 4 16 46 30 0
1 14 5 1 2 14
1 2 3 0 the given system of equations have infinitely many solutions.
2 5 x 2
Example 1.65 [ME-2016 (1 mark)]: The solution to the system of equations
4 3 y 30
is
(a) 6, 2 (b) –6, 2 (c) –6, –2 (d) 6, –2
2 5 x 2
Solution (d): The given system is written as AX B , where A , X ,B .
4 3 y 30
So using Cramer rule, we have x 1 , y 2 , where,
2 5 2 5 2 2
6 20 26 , 6 150 156 and 60 8 52 .
3 30 3 4 30
1 2
4
Thus x 1 156 26 6 and y 2 52 26 2 . Thus option (d) is correct.
Example 1.68 [CH-2016 (2 marks)]: A set of Solution: The given system of linear
simultaneous linear algebraic equations is represented in equations in matrix form can be written as:
a matrix form as shown below. The value (rounded off 4 x4 13 x5 46 …(i);
to the nearest integer) of x3 is _____. 2 x1 5 x2 5 x3 2 x4 10 x5 161 …(ii);
0 0 0 4 13 x1 46 2 x 5 x 3x 61 …(iii);
3 4 5
2 5 5 2 10 x2 161
4 x4 5 x5 30 …(iv);
0 0 2 5 3 x3 61 2 x1 3x2 2 x3 x4 5 x5 81 …(v)
0 0 0 4 5 x4 30 From (i) & (iv), x4 5 and x5 2 ; from
2 3 2 1 5 x5 81 (iii), we get x3 15 .
Matrix Inverse Method: Matrix Inverse method is applicable when number of equations is same
as the number of unknowns, i.e., for (1.5), we have m n . Consider a system of n linear equation in
n unknowns; the system of equations can be written in matrix form Ax B where, the coefficient
matrix A ; the column vectors x and B are given by Eq. 1.6. Multiplying both sides of Ax B by
1
the inverse matrix A , we get A 1 Ax A1 B Ix A1 B x A1 B . So for finding the
solution, we have to find the inverse of the coefficient matrix and the solution given as
1 1 adj ( A)
x A B; A (1.8)
A
Nature of Solution of System of Linear Equations
If A is a non-singular matrix, then the system of equation given by Ax B as a unique solution
given by x A1 B .
If A is a singular matrix, and ( adjA) B 0 , then the system of equation given by Ax B is
consistent with infinitely many solutions.
If A is a singular matrix, and ( adjA) B 0 , then the system of equation given by Ax B is
inconsistent and has no solutions.
Example 1.69 [ME-1997 (1 mark)]: For the following set of simultaneous linear equations:
1.5 x 0.5 y 2 ; 4 x 2 y 3 z 9 ; 7 x y 5 z 10 .
(a) the solution is unique (b) infinitely many solutions exist
(c) the equations are incompatible (d) finite number of multiple solutions exist
1.5 0.5 0 x 2
Solution (a): A 4 2 3 , X y and B 9 are the coefficient matrix and column
7 1 5 z 10
1.5 0.5 0
vectors of the given system. A 4 2 3 1.5(10 3) 0.5(20 21) 10 0 A is non-
7 1 5
singular matrix. Hence the system has unique solution.
Example 1.71 [ME-2012, PI-2012 (2 marks)]: The system of algebraic equation given
x 2 y z 4 ; 2 x y 2 z 5 ; x y z 1 ; has
(a) a unique solution of x 1 , y 1 and (b) only the two solutions of ( x 1 , y 1 and z 1 )
z 1 and ( x 2 , y 1 and z 0 )
(c) infinite number of solutions (d) no feasible solution
1 2 1 x 4
Solution (c): A 2 1 2 , X y and B 5 are the coefficient matrix and column vectors
1 1 1 z 1
1 2 1
1 2
of the given system. A 2 1 2 0 A is singular. Also, C11 3,
1 1
1 1 1
2 2 2 1 2 1 1 1 1 2
C12 0 , C13 3 , C21 3 , C22 0 , C23 3,
1 1 1 1 1 1 1 1 1 1
3 0 3
2 1 1 2 1 2
C31 3, C32 0, C33 3 Cij 3 0 3
1 2 1 2 2 1
3 0 3
3 3 3 3 3 3 4 0
adj ( A) Cij 0
T
0 0 adj ( A) B 0
0 0 5 0 O . Hence the
3 3 3 3 3 3 1 0
system of equation given by Ax B is consistent and has infinitely many solutions.
Gauss Elimination and back substitution: This is a standard elimination method for solving
linear system given by (1.2), where number of equations may or may not be equal to number of
unknowns. To perform Gaussian elimination starting with the system of equations given by Ax B .
a11 a12 a1n | b1
a a22 a2 n | b2
Composing the augmented matrix equation given by Eq. 1.4 as, A .
21
|
am 2 amn | bm
am1
Here, the column vector in the variables x is carried along for labelling the matrix rows. Now
perform elementary row operations to put the augmented matrix into the upper triangular form, i.e.,
a11 a12 a1n | b1 x1 a11 a12 a1n x1 b1
0
a22 a2 n | b2 x2
0
a22 a2 n x2
b
2 (1.9)
|
xn bm
0 0 amn
| bm xn 0 0 amn
Solve the equation of the mth row for xn ; then substitute back into the equation of ( n 1)th row to
obtain a solution for xn 1 ; then substitute back into the equation of (n 2)th row to obtain a solution
for xn 2 and so on. Thus, in Gauss elimination method for solving a system of linear algebraic eqn,
triangularization leads to upper triangular matrix. [This point was asked in ME-1996 (1 mark)].
Nature of Solution of System of Linear a11 a12 a1n | b1
Equations: At the end of Gauss elimination the
form of the coefficient matrix, the augmented c22 c2 n | b2
matrix, and the system itself are called the row |
echelon form; i.e., at the end of Gauss elimination
kr r kr n | br (1.10)
(before the back substitution) the row echelon form
of the augmented matrix is given by Eq. 1.10. Here, | br1
r m , a11 0 , c22 0 , , kr r 0 and all the |
entries in triangle as well as in rectangle are zero. | bm
Now w.r.t solutions of the system with augmented matrix (1.10), there are three possible cases:
Exactly one solution: If r n and br1 , , bm if present, are zero. To get the solution, solve
the nth equation corresponding to (1.10) for xn , then ( n 1)th equation for xn 1 , and so on. up the
line. See Example 1.72.
Infinitely many solution: If r n and br1 , , bm if present, are zero. To obtain any of these
th
solutions, choose values of xr 1 , , xn , arbitrarily. Then solve the r equation for xr , then the
th
(r 1) equation for xr 1 , and so on up the line. See Example 1.73.
No solution: If r m and one of the entries of br1 , , bm is not zero. See Example 1.74.
Example 1.72 [CE-1998 (5 marks)]: Solve the following set of simultaneous equations by Gauss
elimination method. x 2 y z 3 ; x 3z 11 ; 2 y z 1 .
Solution: The augmented matrix of the given system is given as:
1 2 1 | 3 1 2 1 | 3 1 2 1 | 3
1 0 3 | 11 By Applying 0 2 2 | 8 By Applying 0 2 2 | 8
R R R R R R
0 2 1 | 1
2 2 1 0 2 1 | 1
3 3 2 0 0 3 | 9
1 2 1 x 3 x 2 y z 3 (i)
Restoring the transformed matrix equation 0 2 2 y 8 2 y 2 z 8 (ii) .
0 0 3 z 9 3z 9 (iii)
(iii) z 3 ; Putting value of z in (ii), we get, y 1 ; Putting values of z and y in (i), we get,
x 2 . Hence we get a unique solution.
[Similar question was also asked in EE-1995 (5 marks)]
Example 1.74 [CE-2006 (1 mark)]: Solution for the system defined by the set of equations
4 y 3 z 8 ; 2 x z 2 ; 3 x 2 y 5 is
(a) x 0; y 1; z 4 3 (b) x 0; y 1 2; z 2 (c) x 1; y 1 2; z 2 (d) Non-existent
Solution (d): The augmented matrix for the system of linear equations is:
0 4 3 | 8 2 0 1 | 2 2 0 1 | 2
2 0 1 | 2 By Applying 0 4 3 | 8 By Applying 0 4 3 | 8
R2 R1 3 2 R3 3 R1
3 2 0 | 5 R
3 2 0 | 5 0 4 3 | 6
2 0 1 | 2 2 0 1 x 2 2x z 2 (i)
By Applying
0 4 3 | 8 0 4 3 y 8 4 y 3 z 8 (ii) .
R 3 R3 R
2
0 0 0 | 6 0 0 0 z 6 06 (iii)
As (iii) is not valid for any values of x , y and z thus the system has no solution so it is inconsistent.
Example 1.76 [AE-2012 (1 mark)]: The value of k for which the system of equations
x 2 y kz 1 ; 2 x ky 8 z 3 has no solutions is
(a) 0 (b) 2 (c) 4 (d) 8
Solution (c): The augmented matrix for the system of linear equations is
1 2 k | 1 By Applying 1 2 k | 1
2 k 8 | 3 R R 2 R 0 k 4 2( k 4) | 1
2 2 1
Now the system of linear equation have no solution if k 4 0 k 4 .
Example 1.78 [CS-2016 (1 mark)]: Consider the system, each consisting of m linear equations in n
variables.
I. If m n , then all such system have a solution
II. If m n , then none of these systems has a solution
III. If m n , then there exist a system which has a solution.
Which one of the following is CORRECT?
(a) I, II and III are true (b) Only II and III are true (c) Only III is true (d) None of them is true
Solution (c): For statement I: Let we have the system which has 2 equations and 3 unknowns as:
x y z 1 and x y z 1 , which has no solution. So, if m n , then all such system does not
necessary to have a solution. Thus statement I is not correct.
For statement II: Let we have the system which has 3 equations and 2 unknowns as: x y 2 ,
x y 0 and 2 x y 3 , which has a solution ( x 1, y 1 ). So, if m n , then there exist a system
which have a solution. Thus statement II is not correct.
For statement III: Let we have the system which has 2 equations and 2 unknowns as: x y 2 and
x y 0 which has a solution ( x 1, y 1 ). So, if m n , then there exist a system which have a
solution. Thus statement III is correct. Hence option (c) is correct.
Example 1.80 [CH-2007 (1 mark)]: The value of ‘ a ’ for which the following set of equations
y 2 z 0 ; 2 x y z 0 ; ax 2 y 0 have non-trivial solutions, is
(a) 0 (b) 8 (c) –2 (d) 3
0 1 2
Solution (b): 2 1 1 0 1(0 a ) 2(4 a ) 8 a . Also 1 2 3 0
a 2 0
For non-trivial solution 1 2 3 0 0 8 a 0 a 8
[Similar questions were also asked in PI-2010 (1 mark), MN-2013 (2 marks)]
Example 1.81 [MN-2008 (2 marks)]: The solution of the following system of linear equations
x 4 y 3 z 0 ; 3 x 5 y 2 z 0 ; 8 x 10 y 12 z 0 is
(a) (0, 0, 0) (b) (1, 1,1) (c) (2, 1, 2) (d) ( 3, 0,1)
Solution (a): As the given system of linear equation is homogenous with
1 4 3
3 5 2 1(60 20) 4(36 16) 3(30 40) 10 0 . Hence the given system is
8 10 12
consistent with only trivial solutions.
Exercise: 1.2
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
12. If x ( y z ) a , y ( z x) b and z ( x y ) c , where x, y , z are not all zero, then the value of
1 ab bc ca _____.
13. If pqr 0 and the system of equations ( p a) x by cz 0 , ax ( q b) y cz 0 ,
ax by ( r c ) z 0 has a non-trivial solution, then absolute value of {(a p ) (b q) (c r )} is
_____.
14. The system of equations kx y z k 1 , x ky z k 1 , x y kz k 1 , has no solution if
k is
(a) 1 (b) –2 (c) 2 (d) –1
15. The set of equations x y (cos ) z 0 , 3 x y 2 z 0 , (cos ) x y 2 z 0 , where
[0, 2 ) , has non-trivial solutions
(a) for no value of and (b) for all values of and
(c) for all values of and only one value of (d) for all values of and no value of
16. If c 1 and the system of equations x y 1 0 , 2 x y c 0 , bx 3by c 0 is consistent,
then one of the possible real value of b is
(a) 2 (b) 3 (c) 4 (d) 5
17. If a, b, c are in GP with common ratio r1 and p, q, r are in GO with common ratio r2 , and
equations ax py z 0 , bx qy z 0 , cx ry z 0 have only zero solution, then which
one of the following is correct?
(a) r1 1 (b) r2 1 (c) r1 r2 (d) All (a), (b), (c)
18. If a, b, c are non-zero real numbers and if the equations ( a 1) x y z , (b 1) y z x and
(c 1) z x y have a non-trivial solution, then ab bc ca abc _____.
1 0 0
19. Let A 2 1 0 , if x1 and x2 are column matrices such that Ax1 1 0 0 and
T
3 2 1
T
Ax 2 0 1 0 , then x1 x 2
T T T T
(a) 1 1 0 (b) 1 1 1 (c) 1 1 0 (d) 1 1 1
20. The solution of set of equations given by x y z 6 , x y z 2 , 3 x 2 y 4 z 5 is
(a) x 2, y 1, z 3 (b) x 3, y 2, z 1 (c) x 2, y 3, z 1 (d) x 1, y 2, z 3
21. The system of equations x1 2 x2 x3 3 , 2 x1 3 x2 x3 3 , 3 x1 5 x2 2 x3 1 has
(a) infinite number of solutions (b) exactly 3 solutions
(c) a unique solution (d) no solution
22. The system of equations 2 x 6 y 11 0 , 6 x 20 y 6 z 3 0 , 6 y 18 z 1 0 is
(a) consistent (b) inconsistent (c) unique solution (d) infinitely many solutions
23. The system of equations x 2 y 3 z 0 , 3 x 4 y 4 z 0 , 7 x 10 y 12 z 0 has
(a) No non-trivial solution (b) Infinite number of non-trivial solutions
(c) Only trivial solutions (d) An unique non-trivial solution
24. The system of equations 2 x y z 1 , x 2 y z 1 , x y 2 z 2 has
(a) unique solution (b) infinitely many solutions
(c) no solutions (d) a set of only two solutions
25. The system of equations x 2 y 3z 1 , 2 x y 3 z 2 , 5 x 5 y 9 z 4 has
(a) a unique solution (b) infinitely many solutions
(c) no solution (d) exactly 3 solutions
26. The system of equations x y z 0 , 2 x y z 0 , 3 x 2 y 0 has
(a) no solution (b) infinite number of non-trivial solution
(c) only trivial solution (d) trivial and infinite number of non-trivial solutions
A system of axes with this property is called a right-handed system. As shown in Fig. 1.2(a), if the
coordinates of P are ( x, y , z ) then the position vector of r can be written as OP r x i y j z k .
The vectors x i , y j and z k are called the right angled components of r . The scalars x , y , z are
called the components or resolved parts of r in the directions of x axis, y axis and z axis
respectively and ordered triplet ( x, y , z ) is known as coordinates of P whose position vector is r .
The magnitude or modulus or norm of r r x2 y2 z2 .
Direction cosines of r are the cosines of angles that the vector r makes with the positive
direction of x, y , z axes, i.e., cos l x r , cos m y r and cos n z r . Clearly,
2 2 2
l m n 1 . Here POX , POY and POZ and i, j, k are the unit vectors
along OX , OY , OZ axis, respectively.
For 2D – plane as shown in Fig. 1.2(b), i.e., xy , yz , zx the z, x , y respectively, coordinate
becomes zero and we get magnitude and direction cosine. For e.g., as shown in Fig. 1.3, in xy
2 2
plane, we have z 0 , hence, magnitude of r r x y and cos l x r ,
cos m y r and cos n 0 90o .
Addition in component form: If the vectors are defined in terms of i , j and k , i.e., if
u u1 i u2 j u3 k and v v1 i v2 j v3 k , then their sum is defined as
u v (u1 v1 ) i (u2 v2 ) j (u3 v3 ) k . Let we have three vectors u , v and w then we can also
show that the vector addition is associative as well, i.e.,
Associative law: (u v ) w v (u w ) (1.13)
From Eq. 1.12 it follows for each vector u that,
u0 0u (1.14)
Corresponding to u we define a negative inverse ‘ u ’ such that
if u is any nonzero vector; then u is determined uniquely as
shown in Fig. 1.4(a), i.e., it is of same length as u but is directed Figure 1.4: Subtraction of Vector
in the opposite direction (again, u and u have the same length,
the length of u is not negative). For the zero vector we have 0 0 . We denote ‘ u ( v ) ’ as ‘
u v ’ but emphasise that it is really the addition of u and v , as shown in Fig 1.4(b).
Position Vector: If a point O is fixed as the origin in space (or plane) and P is any point, then OP
is called the position vector of P with respect to O . If we say that
P is the point r , then we mean that the position vector of P is r w.r.t. to
some origin O .
AB in terms of the position vectors of points A and B : If a and b
are position vectors of points A and B respectively, as shown in
Fig. 1.5(a). Then, OA a, OB b . In OAB , we have
OA AB OB AB OB OA b a AB (Position vector
of B ) – (Position vector of A ) AB (Position vector of
head) – (Position vector of tail).
Position vector of a dividing point
Internal division: Let A and B be two points with position
vectors a and b respectively, and let C be a point dividing AB
internally in the ratio m : n , as shown in Fig. 1.5(b). Then
m b n a
the position vector of C is given by OC .
mn
External division: Let A and B be two points with position
vectors a and b respectively and let C be a point dividing AB
externally in the ratio m : n , as shown in Fig. 1.5(c). Then
m b n a Figure 1.5: (a) Position Vector (b)
the position vector of C is given by OC . Internal division (c) External division
mn
Scalar Multiplication of a Vector: The scalar multiplication, between a vector u and any scalar
(i.e., a real number) : If 0 and u 0 then u is a vector whose length is times the length
of u and whose direction is same as that of u if 0 , and opposite if 0 ; if 0 and/or u 0
, then u 0 . It follows from this definition that scalar multiplication has the following properties:
( u ) u (1.15) ( )u u u (1.16)
(u v) u v (1.17) 1u u (1.18)
where, , are any scalars and u and v are any vectors.
Example 1.83 [EC-2016 (1 mark)]: Consider the time-varying vector I xˆ 15 cos(t ) yˆ 5sin( t )
in Cartesian coordinates, where 0 is a constant. When the vector magnitude I is at its minimum
value, the angle that I makes with the x axis (in degrees, such that 0 180 ) is _____.
Solution: I {15 cos(t )}2 {5sin( t )}2 225 cos 2 (t ) 25sin 2 (t )
I 200 cos 2 ( t ) 25 100{1 cos(2 t )} 25 125 100 cos(2 t ) , which is minimum
when cos(2t ) 1 2t 180 o t 90 o . We know that t , so the required value of in
o
degrees for which I is minimum is 90 .
u v u v cos (1.19)
In component form, u v (u1 i u2 j u3 k ) (v1 i v2 j v3 k ) u1v1 u2v2 u3v3 (1.20)
Since, i i j j k k 1 and i j j i j k k j k i i k 1 .
If u v 0 is acute angle If u v 0 either u 0 or v 0 or 90o
If u v 0 is obtuse angle [This point was asked in CE-2008 (2 marks)]
Geometrical interpretation of Scalar product: Let u and v be two vectors represented by OA
and OB respectively, as shown in Fig 1.6. Let be the angle between OA and OB . Draw BL OA
and AM OB . From OBL and OAM , we have OL OB cos and OM OA cos . Here OL
and OM are known as projection of v on u and u on v respectively. Hence,
u v u v cos u (OB cos ) u (OL) (Magnitude of u)(Projection of v on u) (1.21)
Also, u v u v cos v (OA cos ) v (ON ) (Magnitude of v)(Projection of u on v ) (1.22)
Thus geometrically interpreted, the scalar product of two vectors is the product of modulus of either
vector and the projection of the other in its direction.
( 3 2) ( 3 2) (1 2) (1 2)
cos 1 cos 1 ( 1 2) 120o
( 3 2) 2 (1 2) 2 2 2
( 3 2) (1 2)
[Similar question was also asked in ME-2004 (1 mark)]
Example 1.85 [IN- Solution: The vector which is perpendicular to the vectors (i j k ) and
2016 (1 mark)]: The i j k
vector that is NOT
perpendicular to the (i 2 j 3k ) is parallel to (i j k ) (i 2 j 3k ) 1 1 1
vectors (i j k ) and 1 2 3
(i 2 j 3k ) is _____. i (3 2) j(3 1) k (2 1) i 2 j k . So the vectors which are parallel
(a) (i 2 j k ) to (i 2 j k ) are given by options (a), (b) and (c). The vector (0i 0 j 0k )
(b) ( i 2 j k ) is parallel to the vector (i 2 j k ) , as the direction of (0i 0 j 0k ) is any
(c) (0i 0 j 0k ) arbitrary direction. So the vector which is not perpendicular to the vectors
(d) (4i 3j 5k ) (i j k ) , (i 2 j 3k ) is given by option (d).
Example 1.86 [MN-2016 (2 marks)]: A force P 2iˆ 5 ˆj 6kˆ acts on a particle. The particle is
moved from point A to point B, where the position vectors of the position vectors of A and B are
6iˆ ˆj 3kˆ and 4iˆ 3 ˆj 2 kˆ respectively. The work done is _____.
Solution: AB (4iˆ 3 ˆj 2kˆ) (6iˆ ˆj 3kˆ ) 2iˆ 4 ˆj kˆ .
So required work done will be W F AB (2iˆ 5 ˆj 6kˆ ) ( 2iˆ 4 ˆj kˆ ) 4 20 6 22
Properties of Scalar Product
The scalar product of two vector is commutative i.e., u v v u .
The scalar product of vectors is distributive over vector addition i.e.,
(a) u ( v w ) u v u w (Left distributivity)
(b) (u v ) w u w v w (Right distributivity)
If m is a scalar and u , v be any two vectors, then ( mu) v m(u v) u ( mv)
If m , n are scalars and u , v be two vectors, then mu nv mn(u v ) ( mnu) v u ( mnv )
For any vectors u and v , we have (i) u ( v ) (u v ) ( u) v ; (ii) ( u) ( v ) u v
u v
Angle between two vectors: If is the angle between u and v , then sin . In terms of
u v
component, if u u1 i u2 j u3 k and v v1 i v2 j v3 k then
(u2 v3 u3v2 ) 2 (u1v3 u3v1 ) 2 (u1v2 u2v1 ) 2
sin (1.27)
(u12 u22 u33 )(v12 v22 v33 )
Right and Left handed system of vectors: Three mutually perpendicular vectors a, b , c form a
right handed system of vector iff a b c ,
b c a , c a b , as shown in Fig. 1.9b.
For e.g., the unit vectors i, j, k form a right-
handed system, i j k , j k i , k i j
, as shown in Fig. 1.9a. The vectors a, b , c ,
mutually perpendicular to one another form
a left handed system of vector iff c b a ,
a c b , b a c , as shown in Fig. 1.9c. Figure 1.9: (a), (b) Right – handed (c) Left – handed system
Example 1.87 [CS-1995 (2 marks)]: A unit vector perpendicular to both the vectors
a iˆ 2 ˆj kˆ and b iˆ ˆj 2kˆ is
(a) (iˆ ˆj kˆ ) 3 (b) (iˆ ˆj kˆ ) 3 (c) (iˆ ˆj kˆ ) 3
(d) (iˆ ˆj kˆ) 3
Solution (a): Since a vector which perpendicular to both the vectors a 2iˆ 2 ˆj kˆ and
b iˆ ˆj 2kˆ is found by taking their cross product, i.e.,
iˆ ˆj kˆ
c a b 1 2 1 iˆ(4 1) ˆj (2 1) kˆ (1 2) 3iˆ 3 ˆj 3kˆ .
1 1 2
3iˆ 3 ˆj 3kˆ 3iˆ 3 ˆj 3kˆ iˆ ˆj kˆ
As we have to find the unit vector, so c c .
2 2 2
3 3 3 3 3 3
Example 1.88 [MN-2007 (2 marks)]: Two sides of a triangle are represented by vectors
a iˆ ˆj kˆ and b iˆ ˆj kˆ . The area (magnitude) of the triangle is
(a) 1 2 (b) 1 (c) 2 (d) 2 2
iˆ ˆj kˆ
Solution (c): Since, a b (iˆ ˆj kˆ) ( iˆ ˆj kˆ ) 1 1 1 2(iˆ ˆj ) So, the area of the
1 1 1
triangle formed by a and b is (1 2) a b (1 2) 2(iˆ ˆj ) 12 12 2
Example 1.89 [AG-2008 (1 mark)]: The cross product of x 2iˆ ˆj and y iˆ 2 ˆj kˆ is
(a) iˆ ˆj 2 kˆ (b) iˆ 2 ˆj 5kˆ (c) iˆ 2 ˆj 5kˆ (d) 2iˆ 4 ˆj
iˆ ˆj kˆ
Solution (c): x y 2 1 0 iˆ(1 0) ˆj (2 0) kˆ (4 1) iˆ 2 ˆj 5kˆ
1 2 1
[Similar question was also asked in AG-2013 (1 mark)]
Example 1.91 [PI-2011 (1 mark)]: If A(0, 4,3) , Solution (a): A vector perpendicular to
B(0, 0, 0) and C (3, 0, 4) are three points defined in both BA 4 ˆj 3kˆ and BC 3iˆ 4kˆ is
x, y , z coordinate system, then which one of the found by taking their cross product, i.e.,
following vectors is perpendicular to both the line vectors
iˆ ˆj kˆ
BA and BC ? BA BC 0 4 3 16iˆ 9 ˆj 12kˆ
ˆ ˆ
(a) 16i 9 j 12 k ˆ ˆ ˆ
(b) 16i 9 j 12 k ˆ
3 0 4
(c) 16iˆ 9 ˆj 12 kˆ (d) 16iˆ 9 ˆj 12 kˆ
Example 1.92
[CE- Solution (a): a b ab sin , where is the angle between a and b .
2011 (2 marks)]: If a
ab
and b are two arbitrary Also, a b ab cos cos sin 1 cos 2
vectors with magnitude ab
a and b , respectively, 2
2 ab a 2 b 2 ( a b) 2
a b will be equal to sin 1
ab ab
2 2
2
(a) a b (a b) a 2b 2 ( a b ) 2
a b ab sin ab a 2b 2 ( a b ) 2
(b) ab a b ab
2 2
2 2
(c) a b (a b) a b a 2b 2 ( a b) 2
(d) ab a b
Example 1.95 [MN-2010 (2 marks)]: The volume of tetrahedron with vertices at (0, 0, 0), (1, 0, 0),
(0, 1, 0) and (0, 0, 1) is
(a) 1/2 (b) 1/4 (c) 1/6 (d) 1/8
Solution (c): Let a 0i 0 j 0k , b i 0 j 0k , c 0i j 0 k and d 0iˆ 0 ˆj kˆ ; so volume
ˆ ˆ ˆ ˆ ˆ ˆ ˆ ˆ ˆ
of tetrahedron is given as,
(1 6) b a ca
d a (1 6) (b a) (c a ) (d a ) (1 6) iˆ ( ˆj kˆ ) (1 6)(iˆ iˆ) 1 6
Specifically, let the point O , as shown in Fig. 1.12, be the origin of a Cartesian coordinate system
with the usual base vectors î , ĵ and k̂ . And let the system be oriented so that î is aligned with b
and k̂ coincides with n̂ , as shown in Fig. 1.12. Then ĵ lust lie in the b , c plane so we can express
b b1 ˆi , c c1 ˆi c2 ˆj and a a1 ˆi a2 ˆj a3 kˆ . Hence b c b1c2 kˆ and
a b c (a iˆ a ˆj a kˆ ) (b c kˆ )
1 2 3 1 2
a b c a1b1c2 ˆj a2b1c2 ˆi (a1c1 a2 c2 )b1 iˆ a1b1 (c1 ˆi c2 ˆj) (a1c1 a2c2 ) v a1b1 w
But a1c1 a2 c2 a c and a1b1 a b . Hence a (b c) (a c) b (a b )c
Example 1.96 [AE-2008 (2 marks)]: Which of the following is true for all choices of vectors p , q ,
r?
(a) p q q r r p 0 (b) ( p q ) r (q r ) p ( r p ) q 0
(c) p ( q r ) q ( r p ) r ( p q ) 0 (d) p (q r ) q (r p ) r ( p q ) 0
Solution (d): Using vector triple product formula a (b c) (a c) b (a b ) c ;
Option (d) can be written as,
p ( q r ) q (r p ) r ( p q ) ( p r ) q ( p q ) r (q p ) r (q r ) p ( r q ) p ( r p ) q 0 .
Exercise: 1.3
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
(a) 1.5 i 1.5 j 4 k (b) 1.5 i 1.5 j 4 k (c) 1.5 i 1.5 j (d) None of these
19. A particle is acted upon by constant forces 4 i j 3 k and 3 i j k which displace it from a
point i 2 j 3 k to the point 5 i 4 j k . The work done in standard units by the force is given
by _____.
2 2 2 2
20. If a is any vector, then (a i) (a j) (a k ) k a , where k _____
21. The sine of the angle between the vectors a 3 i j k , b 2 i 2 j k is
(a) 74 99 (b) 25 99 (c) 37 99 (d) 5 41
22. The vectors c , a x i y j z k and b j are such that a, c, b form a right handed system, then
c is
(a) z i x k (b) 0 (c) y j (d) z i x k
23. If a i j k , b i 3 j 5 k and c 7 i 9 j 11 k , then the area of the parallelogram having
diagonals a b and b c is
(a) 4 6 (b) 4 3 (c) 3 6 (d) 3 3
24. The position vectors of the vertices of a quadrilateral ABCD are a, b, c, d , respectively. Area of
the quadrilateral formed by joining the middle points of its sides is
(a) (1 4) a b b d d a (b) (1 4) b c c d a d b a
(c) (1 4) a b b c c d d a (d) (1 4) b c c d d b
25. Three forces i 2 j 3 k , 2 i 3 j 4 k and i j k are acting on a particle at the point (0,1, 2) .
The magnitude of the moment of the forces about the point (1, 2, 0) is
(a) 2 35 (b) 6 10 (c) 4 17 (d) None of these
26. If u, v, w are three non-coplanar vectors, then (u v w ) [(u v ) ( v w )] equals
(a) 0 (b) u ( v w ) (c) u ( w v ) (d) 3u ( v w )
27. The value of ‘ a ’ so that the volume of parallelopiped formed by i a j k ; j a k and a i k
becomes minimum is
(a) 3 (b) –3 (c) 3 (d) 1 3
28. Let a, b and c be non-zero vectors such that (a b ) c (1 3) b c a . If is the acute angle
between the vectors b and c , then sin equals
(a) 2 2 3 (b) 2 3 (c) 2 3 (d) 1 3
29. Let the vectors a, b, c, d be such that (a b) (c d) 0 . Let P1 and P2 be planes determined by
pair of vectors a, b and c, d respectively. Then the angle (in degrees) between P1 and P2 is ____.
30. If a i j k , a b 1 and a b j k , then b
(a) i (b) i j k (c) 2 j k (d) 2i
31. If be the angle between the unit vectors a and b, then cos( 2)
1 1 a b a b
(a) a b (b) ab (c) (d)
2 2 a b a b
32. If a and b are unit vectors and (a b) is also a unit vector, then the angle (in degrees) between
a and b is _____.
33. If F1 i j k , F2 i 2 j k , F3 j k , A 4 i 3 j 2 k and B 6 i j 3 k , then the
scalar product of F1 F2 F3 and AB will be _____.
34. The horizontal force and the force inclined at an angle 60o with the vertical, whose resultant is in
vertical direction of P kg, are
(a) P, 2 P (b) P, P 3 (c) 2 P, P 3 (d) None of these
Given a vector space V over a field K , we refer to the elements of the field K as scalars (real or
complex). We will restrict our self to real vector space, hence, we use field R instead of field K .
Given an element v of the vector space V , we shall refer to elements of V that are of the form c v
for some scalar c as scalar multiples of v . A vector space over the field R consists of a set V on
which is defined an operation of addition (denoted by ), associating to elements u and v of V an
element u v of V ; and an operation of multiplication by scalars, associating to each element c of
R and to each element v of V an element c v of V , where the following properties holds true:
u, v V u v V
u v v u u , v V (i.e., vector addition is commutative);
(u v ) w u ( v w ) u, v, w V (i.e., vector addition is associative);
an element 0 of V (known as the zero element) with the property that u 0 u u V ;
given any element u V , there exists an element u V with the property that u ( u) 0
c (u v ) c u c v u, v V and c R ; (c d )u (c u d u) u V and c, d R
c ( d u) (cd ) u u V and c, d R ;
1u u u V , where, 1 is the multiplicative identity element of the field R .
For justifying the above axioms, let, we have 2 - tuple vector (i , j ) , here this 2 - tuple vector is
making a ‘two dimensional vector space V ,’ which can be considered as a plane, which contains
infinite vectors (in the form of i j , , R ). Now, justify yourself, by putting u i & v j .
A vector space V over a field R is said to be trivial if it consists of a single element (which must
then be the zero element of V ). A vector space with more than one element is said to be non-trivial.
Let u , v be elements of a vector space V , then a unique element x of V satisfying x v u
Let V be a vector space over a field R . Then c 0 0 and 0 v 0 c R , v V .
Let V be a vector space over a field R . Then ( c ) v (c v ) , c ( v ) (c v ) c R , v V .
Let V be a vector space over a field R . Then u v u ( 1) v u, v V .
Let V be a vector space over a field R . Let c v 0 , then either v 0 or c 0 c R , v V .
A vector space is free of any dependence of a ‘coordinate system.’
Let we have two ‘2 - tuple’ vectors, a i j & b 3i 2 j . Solving these two vectors in terms of
i & j , i 2a b & j 3a b . Let there is another vector, v 4i 5 j , then in terms of a & b ,
we get, v 7 a b . Thus, v (4, 5) w.r.t. ( i , j ); v (7, 1) w.r.t. ( a , b ) . Hence, the vector v is
same but the two ‘2 - tuple’ are not same (as shown in Figure). So, we can conclude that, v ci d j
has unique representation w.r.t. ( i , j ) coordinate system; v a b has unique representation
w.r.t. ( a , b ) coordinate system. Hence, in general, v exist independently of choice of (i , j ) or ( a , b ) .
option (d) is not correct. Now we have to decide between option (b) or (d). Since, the vectors are
T
linearly independent if there is a trivial solution of v1 v 2 v3 x 0 , where x c1 c2 c3 .
The system of equation and its augmented matrix is given as:
1 2 5 c1 0 1 2 5 | 0 By Applying 1 2 5 | 0
1 3 6 c 0 and 1 3 6 | 0 R R R 0 1 1 | 0
2 2 1
2
1 1 4 c3 0 1 1 4 | 0 R 3 R3 R
2 2 |
1 0 0
By Applying 1 0 3 | 0 1 0 3 c1 0 c1 3c3 0 (i)
R1 R1 2 R2 0 1 1 | 0 0 1 1 c2 0 c2 c3 0 (ii) (after restoring
R R 2
3 3 R
2 0 0 0 | 0
0 0 0 c3 0 00 (iii)
the transformed matrix equation). From (i) and (ii) c1 3c2 3c3 and so we can say that all the
scalars c1 , c2 and c3 are not zero. Thus the given vectors are linearly dependent.
Bases for vector space : A finite set S of vector space V is called a basis for V provided that
(1) The vectors in S are linearly independent (2) The vectors in S span V , i.e. span (S) V . For e.g.,
the set 1 0
T T
0 , 0 1 0 , 0 0 1
T
is a basis for R .3
Example 1.103 [EC-2014 (1 mark)]: For matrices of same dimension M , N and scalar c , which
one of these properties DOES NOT ALWAYS hold?
(a) ( M T )T M (b) (cM )T c ( M )T (c) ( M N )T M T N T (d) MN NM
Solution (d): Since matrix multiplication is not commutative in general; so for matrices of same
dimension M and N , MN NM does not always holds.
Similarly, the dot or scalar product of u (u1 , u2 , , un ) and v (v1 , v2 , , vn ) is defined as,
n
u v u1v1 u2 v2 un vn ui vi (1.29)
i 1
Orthogonality of two vectors: If u (u1 , u2 , , un ) and v (v1 , v2 , , vn ) are two non – zero
vectors such that u v 0 then the angle between the two vectors u and v becomes 2 then we
say that the two vectors u and v orthogonal. If we have a set of non – zero vectors {u1 , u 2 , , u k }
then we have orthogonal set of vectors if and only if every vector in the set is orthogonal to every
other one, i.e., u i u j 0 i j .
Two vectors v1 and v 2 are orthogonal iff each is non-zero and the dot product v1 v 2 0 , i.e.,
( v1 )T ( v 2 ) 0 . For e.g., the vectors 3 2 and 2 3 are orthogonal because
T T
3 2 2 3 6 6 0 ; the vectors 3 2 and 1 3 are not orthogonal because
T
3 2 1 3 3 6 0 .
Three vectors v1 , v 2 and v 3 are orthogonal iff each is non-zero and they are pairwise
orthogonal, i.e., ( v1 )T ( v 2 ) 0; ( v 2 )T ( v 3 ) 0; ( v 3 )T ( v1 ) 0 . For e.g., the vectors
T
2 0 0 , 0 2 0 , 0 0 2 are orthogonal because, 2 0 0 0 2 0 0;
T T
0 2 0 0 0 2 0; 0 0 2 2 0 0 0 .
n
If n vectors ( v1 , v 2 , , v n ) each of which is in R are orthogonal, then they are surely linearly
n n
independent and hence span R and therefore form a basis for R . For, e.g., The vectors
2 0 0 , 0 2 0 , 0 0 2 are orthogonal and hence are linearly independent and hence
3 3
span R . They form a basis for R .
Orthogonal Projection: Let u and v be two vectors (not orthogonal) in a dot product space V .
uv
If v 0 , then orthogonal projection (OP) of u onto v is given by (OP) v u v
vv
uv
The vector ( w ) which is orthogonal to v is given as: w u (OP) v u u v
vv
The orthogonal projection (OP) of w onto the subspace spanned by the orthogonal vectors u
w u w v
and v is given by w (OP)u w (OP) v w u v .
uu v v
Normalization: A vector of unit length is called a unit vector. Any non – zero vector u can be
scaled to have unit length by multiplying it by 1 u so we say that the vector û as, uˆ u u .
Orthonormal: A set of vectors {u1 , u 2 , , u k } is said to be orthonormal if it is orthogonal, i.e.,
u i u j 0 i j and if each vector is normalized, i.e., u j u j 1 j . [This point was asked
Example 1.105 [ME-1995 (2 mark)]: Among the following, the pair of vectors orthogonal to each
other is
(a) [3, 4, 7],[3, 4, 7] (b) [1, 0, 0],[1,1, 0] (c) [1, 0, 2],[0, 5, 0] (d) [1,1,1],[ 1, 1, 1]
Solution (c): The given two vectors are orthogonal if their dot product is zero. We can verify that the
dot product of the vectors is zero if we take the vectors in option (c).
Example 1.107 [EE-2006 (2 marks)]: The following vectors is linearly dependent upon the solution
to the previous problem
T T T T
(a) 8 9 3 (b) 2 17 30 (c) 4 4 5 (d) 13 2 3
T
Solution (b): We have to find which of the given vectors is linearly dependent with 6 3 6
T
and 4 2 3 . For linearly dependent vectors the determinant of the matrix, whose first two
T T
column elements are 6 3 6 , 4 2 3 and third column element is taken one of the
given vectors in options, must be zero. By doing so we find option (b) is correct as
6 4 2
3 2 17 6( 60 51) 4( 90 102) 2( 9 12) 0
6 3 30
Example 1.108 [CS-2007 (2 marks)]: Consider the set of (column) vectors defined by
T
T
X x R 3 x1 x2 x3 0 , where x x1 , x2 , x3 . Which of the following is TRUE?
(b) 1, 1, 0 , 1, 0, 1 is linearly independent set, but it does not span
T T
X and therefore is not a
basis of X
3
(c) X is not a subspace of R
(d) None of these
Solution (a): To be basis for subspace X , two conditions are to be satisfied: (i) the vectors have to be
linearly independent (ii) they must span X . Here
X x R 3 x1 x2 x3 0 , where
T
xT x1 , x2 , x3 . Since, 1, 1, 0
T
, 1, 0, 1
T
is a linearly independent set because one cannot be
obtained from another by scalar multiplication. The fact that it is independent can also be established
1 1 0
by seeing the rank of is 2. Next, we need to check if the set spans X . Here
1 0 1
k1 k 2
X x R x1 x2 x3 0 , the general infinite solution of X k1 . Choosing k1 , k2 as
3
k1
k1 0 k1 k T
k k and k 0 , we get two linearly independent solutions for X as, X k 0 k
2 2
T T
or X k k 0 . Now since both of these can be generated by linear combination of 1, 1, 0
T
and 1, 0, 1 , the set spans X . Since we have shown that the set is not only linearly independent but
also spans X , therefore by definition it is a basis for the subspace X .
Example 1.110 [EE-2007 (2 marks)]: Let x and y be two vectors in a 3 dimensional space and
x, x x, y
x, y denote their dot product. Then the determinant det
y, x y , y
Example 1.112 [AE-2013 Solution (d): From the given matrix we have three vectors
(2 marks)]: Values of a , T T
u1 1 3 1 3 1 3 , u 2 1 2 0 1 2 and
b and c , which render the
matrix T
u3 a b c . For these set of vectors to be orthonormal if it is
1 / 3 1 / 2 a orthogonal and if each vector is normalized. For the set of vectors to be
Q 1 / 3 0 b orthogonal,
u i u j 0 i j u1 u 2 0 , u 2 u 3 0 and u 3 u1 0 . As
1 / 3 1 / 2 c u u 0 is true; and
1 2
orthonormal are,
respectively a c
u 2 u3 0 0ac …(i)
1 1 2 2
(a) , ,0
2 2 a b c
u 3 u1 0 u 2 u 3 0 0 a b c 0 …(ii)
1 2 1 3 3 3
(b) , , For each vector to be normalised, the given vectors must have unit
6 6 6
modulus, i.e., u1 u1 1 , u 2 u 2 1 and u 3 u3 1 . As u1 u1 1 ,
1 1 1
(c) , ,
3 3 3 u 2 u 2 1 ; and u 3 u 3 1 a 2 b 2 c 2 1 …(iii)
2
1 2 1 Putting (i) and (ii) in (iii) we get a 1 6 a 1 6 and (i)
(d) , ,
6 6 6 c 1 6 and (ii) b 1 6 .
So from the given options we can say that option (d) is correct.
Example 1.113 [MN-2014 (1 mark)]: The value of k for which the vectors a 2iˆ 3 ˆj and
b kiˆ 4 ˆj are orthogonal to each other is _____.
Solution: The given two vectors are orthogonal if their dot product is zero, i.e. a b 0
(2iˆ 3 ˆj ) ( kiˆ 4 ˆj ) 0 2k 12 0 k 6 . So the answer is k 6 .
Example 1.115 [EC-2016 (2 marks)]: If the vectors e1 (1, 0, 2) , e 2 (0,1, 0) and e3 ( 2, 0,1)
3
from an orthogonal basis of the three-dimensional real space , then the vector u (4, 3, 3) 3
can be expressed as
(a) u (2 5)e1 3e 2 (11 5)e3 (b) u (2 5)e1 3e 2 (11 5)e3
(c) u (2 5)e1 3e 2 (11 5)e3 (d) u (2 5)e1 3e 2 (11 5)e3
Solution: The given vectors e1 , e 2 , e3 form an orthogonal basis of the three-dimensional real space
3 . So any vector u 3 can be expressed as u x1e1 x2 e 2 x3e3 . Thus we have
(4, 3, 3) x1 (1, 0, 2) x2 (0,1, 0) x3 ( 2, 0,1) (4,3, 3) ( x1 2 x3 ), ( x2 ), (2 x1 x3 )
( x1 2 x3 ) 4 …(i), x2 3 …(ii), (2 x1 x3 ) 3 …(iii)
(i) 2 (iii) 5 x1 2 x1 2 5 ; from (iii), we have x3 3 2( 2 5) 3 (4 5) 11 5 .
So we have, u ( 2 5)e1 (3)e 2 ( 11 5)e 3 . Hence option (d) is correct.
Rotation Matrix
Computing the rotation matrix when axis and angle are given: Given the unit vector axis a
and angle , the rotation matrix is obtained by applying the following relation,
R I 2 sin 2 ( 2) aaT I sin A , where I is the identity matrix, aaT denotes a matrix
corresponding to a 3 3 Cartesian matrix with components (aa )ij ai a j , written as
a1a1 a1a2 a1a3 0 a3 a2
[aa ]ij a2 a1
T
a2 a 2 a2 a3 . The matrix [ A ] a3
0 a1 is the skew symmetric
a3 a1a3 a2 a3 a3 a2 a1 0
matrix, associated with the axis a .
Computing axis and angle when the rotation matrix is given: Given R , below an algorithm
for finding the principal axis a and principal angle of rotation:
Step 1: Compute c (1 2)(tr ( R ) 1) (1 2)( R11 R22 R33 1)
Step 2: Decide whether the angles is an integral multiple of 180 o by checking the cosine and
follow one of the following:
If c 1 , proceed to step 3
c 1 , then R I . Hence, the angle of rotation is 0 and the axis of rotation is arbitrary, so
we can set it to {1, 0, 0} or any other convenient unit vector. Go to step 6.
If c 1 , then the angle of rotation is 180 o . To find the axis of rotation, simply
normalize any nonzero column of R I . Go to step 6.
Step 3: Compute the angle of rotation cos 1 c , 0 180 o
Step 4: Compute sin 1 cos 2
R32 R23 R13 R31 R21 R12
Step 5: Compute the axis of rotation a1 , a2 , a3
2 sin 2 sin 2 sin
Step 6: Stop
0 0 1 Solution (d):
Example 1.116 [IN-2009 (2 marks)]: The matrix P 1 0 0 c (1 2) tr ( P) 1
0 1 0 c (1 2)(0 0 0 1)
c 1 2
rotates a vector about the axis 1 1 1 by an angle of
(a) 30o (b) 60o (c) 90o cos 1 (1 2) 120o
Echelon Form of a Matrix: A matrix A is said to be in Echelon form if either A is null matrix or A
satisfies the following conditions: (i) Every non-zero row in A precedes every zero row; (ii) The
number of zeros before the first non-zero element in a row is less than the number of such zeros in the
next row. It can be easily proved that the rank of a matrix in Echelon form is equal to the number of
non-zero rows (or linearly independent rows) of the matrix. For e.g., consider the two matrices
0 3 2 1 0 2 5
A 0 0 2 5 and B 0 0 1 . The rank of the matrix A is r ( A) 2 because it is in
0 0 0 0 0 0 4
Echelon form and it has two non-zero rows; but the matrix B is not in Echelon form, because the
number of zeros in the second row is not less than the number of zeros in the third row. By applying,
0 2 5
R3 R3 4 R2 , matrix B B 0 0 1 , which is in Echelon form and contains two non-zero
0 0 0
rows and hence r ( B ) 2 . It is to be noted that ‘Elementary transformations do not alter the rank
of matrix.’
Algorithm for finding the rank of a matrix: Let A aij mn be an m n matrix. The
following are the steps for finding the number of non-zero row (or number of independent rows)
which is equal to the rank of matrix.
Step 1: Using elementary row transformations make a11 1 .
Step 2: Make a21 , a31 , , am1 all zeros by using elementary transformation: R2 R2 a21 R1 ,
R3 R3 a31 R1 , , Rm Rm am1 R1 .
Step 3: Make a22 1 by using elementary row transformations.
Step 4: Make a32 , a42 , , am 2 all zeros by using elementary transformation: R3 R3 a32 R2 ,
R4 R4 a42 R2 , , Rm Rm am 2 R2 .
The process used in Step 3 and 4 is repeated utpo ( m 1) th row. Finally we obtain a matrix in Echelon
form which is equivalent to matrix A . The rank of A will be equal to number of non-zero rows in it.
0 0 3
Example 1.117 [CS-1994 (1 mark)]: The rank of the matrix A 9 3 5 is
3 1 1
(a) 0 (b) 1 (c) 2 (d) 3
Solution: For finding the rank of matrix, we have to find the number of linearly independent row
3 1 1 1 1 3 1 3
By Applying By Applying
vectors. A 9 3 5 A 9 3 5
R R R1 R1 3
3 1
0 0 3 0 0 3
1 1 3 1 3 By Applying 1 1 3 1 3
By Applying
A 0 0 2 3 A 0 0 2
R2 R2 9 R1
R3 R3 R2
0 0 3
2 0 0 0
As we have two independent rows, so its rank is 2.
[Similar questions were also asked in ME-1994, CS-2002 (1 mark), EC-2006 (2 marks)]
Example 1.118 [EE-1994 (1 mark)]: A 5 7 matrix has all its entries equal to –1. The rank of the
matrix is
(a) 7 (b) 5 (c) 1 (d) 0
Solution (c): In the given matrix if we choose operation Ci Ci C1 , i 2, 3, 4, 5 ; then we get C1
containing all elements equal to –1 and all other columns C i ( i 2, 3, 4, 5 ) containing all elements
equal to 0. Hence we have only one non-zero column or independent column, so its rank is 1.
Example 1.119 [CS-1995 (1 mark)]: The rank of the following ( n 1) (n 1) Solution (a): If we
carefully notice
1 a a2 an
the given matrix,
1 a a2 an we found we have
matrix A , where a is a real number is only one
independent row
as all the elements
1 a a a
2 n in each row are
same. So its rank
(a) 1 (b) 2 (c) n (d) Depends on the value of a
is 1.
4 2 1 3 (a) 4
(b) 3
Example 1.120 [CE-2003 (1 mark)]: The rank of the matrix 6 3 4 7 is (c) 2
2 1 0 1 (d) 1
Solution (c): For finding the rank of matrix, we have to find the number of linearly independent row
By Applying 1 1 2 1 4 3 4 By Applying 1 1 2 1 4 34
vectors. 6 3 4 7 R2 R2 6 R1 0 0 5 2 5 2
R
R1 1
4 2 1 0 1 R R
3 3 2 R
1 0 0 1 2 1 2
By Applying 1 1 2 1 4 3 4
R 0 0 5 2 5 2 . We have two independent row vectors, so its rank is 2.
R3 R3 2
5 0 0
0 0
[Similar question was also asked in CE-2014 (2 marks)].
T
Example 1.121 [EE-2007 (1 mark)]: x x1 x2 xn is an n tuple non-zero vector. The
n n matrix V xxT
(a) has rank zero (b) has rank 1 (c) is orthogonal (d) has rank n
Solution: After multiplication, We have,
x12 x1 x2 x1 x3 x1 xn x12 x1 x2 x1 x3 x1 xn
2 By Applying
x2 x1 x2 x2 x3 x2 xn 0 0 0 0
2 Ri R1
V x3 x1 x3 x2 x3 x3 xn Ri V 0 0 0 0 . As
xi x1
for all 1 i n
x x xn x2 xn x3 2
xn 0 0 0 0
n 1
Example 1.122 [PI-2007 (2 marks)]: q1 , , qm are n dimensional vectors, with m n . This set of
vectors is linearly dependent. Q is the matrix with q1 , , qm as the columns. The rank of Q is
(a) less than m (b) m (c) between m and n (d) n
Solution (a): As Q is a m n matrix so rank (Q ) min( m, n) rank (Q ) m , as m n . But the
column vectors are linearly dependent so rank (Q) must not be m rank (Q ) m .
Rank of a Matrix: A number r is said to be the rank of an m n matrix A if (i) Every square sub-
matrix of order ( r 1) or more is singular; (ii) There exists at least one square sub-matrix of order r
which is non-singular. Thus, the rank of a matrix is the order of the highest order non-singular
square sub-matrix. For e.g., if A is a given matrix of rank 2, then every square sub-matrix of order 3
or more is singular and there exists at least one square sub-matrix of order 2 which is non-singular.
If A is a non-singular square matrix of order n , then its rank is n .
If A is a m n matrix, then r ( A) min(m, n) . Because the largest order square sub-matrix of A
will be of the order equal to min( m, n) .
1 4 8 7
0 0 3 0
Example 1.123 [CS-1998 (2 marks)]: The rank of the matrix A is
4 2 3 1
3 12 24 2
(a) 3 (b) 1 (c) 2 (d) 4
Solution (d): We have 4 4 matrix. For finding its determinant we will expand it along second row
1 4 7
as it contains maximum number of zeros. A 3 4 2 1 3 266 798 0 A is a non-
3 12 2
singular square matrix of order 4, so its rank is 4.
3 2 12
Example 1.124 [ME-1999 (1 mark)]: Rank of the matrix is A 6 4 8
12 8 36
(a) 1 (b) 2 (c) 3 (d) 2
Solution (b): We first check the determinant of the given matrix have 3 3 matrix.
3 2 12 3 2 12
A 6 4 8 2 4 3 2 4 (by taking common –2 and 3 from R2 and R3 respectively).
12 8 36 3 2 9
As we have C1 and C 2 are same, so A 0 A is a singular matrix and so its rank is not 3. Now let
us choose the square sub-matrix of order 2 2 by leaving first column and third row we get
2 12
16 48 64 0 which is non-singular. As we get a non-singular square sub-matrix of
4 8
order 2 2 , hence the given matrix has rank 2.
[Similar question was also asked in BT-2012 (2 marks)]
n
The row space Row( A ) of A is the subspace of R spanned by the rows of A . The
dimension of Row( A ) is called row rank of A .
m
The column space Col( A ) of A is the subspace of R spanned by column of A . The
dimension of Col( A ) is called the column rank of A .
If A is m n matrix; and B is n k matrix, then
rank( AB ) rank( A) rank( AB ) rank( B )
rank( AB ) min rank( A), rank( B ) rank( AB ) rank( A) rank( B ) n
If A and B are both n n matrices with rank n , then then rank( AB ) n
Example 1.125 [ME-2001 (2 marks)]: The rank of a 3 3 matrix C ( AB) , found by multiplying a
non-zero column matrix A of size 3 1 and a non-zero row matrix B of size 1 3 , is
(a) 0 (b) 1 (c) 2 (d) 3
Solution (b):
a11 a11 a11b11 a11b12 a11b13
Let A a21 and B b11 b12 b13 , then C a21 b11
b12 b13 a21b11 a21b12 a21b13
a31 a31 a31b11 a31b12 a31b13
By Applying a11b11 a11b12 a11b13
R2 R2 a21 R1 a11 C 0 0 0 we have one independent row, so rank is 1.
R R 3 a31 R1
a
3 11
0 0 0
0 1 1
Example 1.129 [IN-2013 (1 mark)]: The dimension of the null space of the matrix 1 1 0 is
1 0 1
(a) 0 (b) 1 (c) 2 (d) 3
Solution (b): For finding the dimension of the null space of the matrix, first we have to find the
dimension of the matrix; and for the dimension of any matrix we have to find its rank.
1 1 0 1 1 0 1 1 0
By Applying By Applying By Applying
0 1 1 0 1 1 0 1 1
R R R R R R R R
1 2 1 0 1
3 3 1
3 3 2 0 0 0
0 1 1
So we have two independent row vectors so its rank is 2 and thus the given matrix is of 2 dimensional.
Hence the dimension of the null space of the given matrix is 3 2 1 .
p q
Example 1.130 [EE-2014 (1 mark)]: Two matrices A and B are given as: A ;
r s
p2 q2 pr qs
B . If the rank of matrix A is N , then the rank of matrix B is
pr qs r 2 s2
(a) N 2 (b) N 1 (c) N (d) 2N
p q By Applying p2 q 2 q
Solution (c): The given matrix A A
r s C 1 pC1 qC
2 pr qs s
By Applying p 2 q 2 pr qs
A 2
B . As matrix B is obtained by
C2 ( r p )C1 ( ps rq) p C 2 2
pr qs r s
elementary row transformation from matrix A , so if rank ( A) N rank ( B ) N .
Example 1.131 [EE-2016 (2 marks)]: Let A be a 4 3 real matrix with rank 2. Which one of the
following statement is TRUE?
T T
(a) Rank of A A is less than 2 (b) Rank of A A is 2
T T
(c) Rank of A A is greater than 2 (d) Rank of A A can be any number between 1 and 3
Solution (b): rank [ A43 ] 2 [ A3T4 ] 2 . As rank( A B ) min{rank( A), rank( B)} .
T T
Order of A A is 4 4 ; and order of AA is 3 3 .
So rank( AT A) min{rank( A), rank( AT )} rank( AT A) min(2, 2) rank( AT A) 2
Similarly, rank( AAT ) min{rank( A), rank( AT )} rank( AAT ) min(2, 2) rank( AAT ) 2 .
Thus from given options, option (b) is correct.
Number of Solutions for non-homogeneous Linear System in terms of rank: Consider the
system of linear equation given by Eq. 1.2 and in matrix notation Eq. 1.2 can be written as Ax B ,
where A a jk O is the coefficient matrix; x x1k 1n and B b1k 1m are the column
m n
Example 1.132 [EE-1998 (1 mark)]: A set of linear equation is represented by the matrix equation
AX B . The necessary condition for the existence of a solution for this system is:
(a) B must be linearly dependent on the columns of A (b) A must be invertible
(c) B must be linearly independent on the columns of A (d) None of the above
Solution (a): For having a solution for the matrix equation AX B , the rank ( A) = rank ( A | B) n .
If rank ( A) n , then A is invertible. If rank ( A) n then A is not invertible. Hence option (b) is
not correct. If rank ( A | B) n then the column of B may or may not be linearly dependent on the
columns of A . If rank ( A | B) n then the column of B must be linearly dependent on the columns
of A . Hence option (a) is correct and (c) is not correct.
Example 1.133 [CS-2003 (2 marks)]: Consider the following system of linear equations:
2 1 4 x
4 3 12 y 5 . Notice that the second and the third columns of the coefficient matrix are
1 2 8 z 7
linearly dependent. For how many values of , does this system of equations have infinitely many
solutions?
(a) 0 (b) 1 (c) 2 (d) Infinitely many
4
2 1 4 |
2 1
Solution (b): Given: A 4 3 12 , A | B 4 3 12 | 5 . We have infinitely many
1 2 8 1 2 8 | 7
solutions for the system given if rank ( A ) rank ( A | B ) 3 (as we have 3 unknowns).
Since rank ( A ) is 2 as we have two linearly independent column.
The rank ( A | B ) is evaluated as:
1 1 2 2 | 2 By Applying 1 1 2 2 | 2
By Applying
4 3 12 | 5 R2 R2 4 R1 0 1 4 | 5 2
R1 R1 2
1 2 8 | 7 3 2 6 | 7 2
R R R
3 3 1 0
By Applying 1 1 2 2 | 2
0 1 4 | 5 2 For rank ( A | B ) rank ( A ) 2, we must
3
R3 R3 R1
2 0
0 0 | 11( 2) 7 2
have 11( 2) (7 2) 0 7 11 . Hence we have one value of .
Example 1.134 [IN-2007 (2 marks)]: Let A aij ,1 i, j n , with n 3 and aij i j . Then the
rank of A is
(a) 0 (b) 1 (c) n 1 (d) n
1 2 3 n 1 2 3 n
2 4 6 2n 0 0 0 0 By Applying
Solution (b): Let A 3 6 9 3n A 0 0 0 0 Rk Rk 2 Rk 1
2 k n
n 2n 3n n
2
0 0 0 0
As we have only one independent row hence its rank is 1.
Solution (b): From the given data we can say that A is m n matrix. m n rank( A) m . If
rank( A) rank( A | b) m n The system have infinitely many solutions.
Example 1.136 [EE-2008 (1 mark)]: If the rank of a (5 6) matrix Q is 4, then which one of the
following statements is correct?
(a) Q will have four linearly independent rows and four linearly independent columns
(b) Q will have four linearly independent rows and five linearly independent columns
(c) QQT will be invertible (d) QT Q will be invertible
Solution (a): The rank of a (5 6) matrix Q is 4 Q has four linearly independent rows and four
linearly independent columns. Hence option (a) is correct and (c) is wrong. Also, as Q has rank
4 min(5, 6) so Q 0 Q 1 does not exist as Q is singular matrix; so (QQT ) 1 (QT ) 1 (Q ) 1
and (QT Q) 1 Q 1 (QT ) 1 both does not exists. So options (b) and (d) both are wrong.
1 1 1 | 3 By Applying 1 1 1 | 3
By Applying 4 0 7 | 1 R2 R2 4 R1 0 4 3 | 11
For r ( A | B ) :
3 2
R1 R 3 0 | 1 R3 R3 3R1 0 1 3 | 8
1 2 7 | 0 R4 R4 R1 0
3 6 | 3
1 1 1 | 3 1 1 1 | 3
By Applying
By Applying 0 1 3 4 | 11 4 0 1 34 | 11 4
R3 R3 R2
R R 4
2 2 0 1 3 | 8
0 0 15 4 | 21 4
R R 3 R
0 0
4 4 2
0 3 6 | 3 15 4 | 21 4
1 1 1 | 3
By Applying 0 1 3 4 | 11 4
. As we have 3 non-zero rows so r ( A | B ) 3 .
R R R
4 4 3 0 0 15 4 | 21 4
0 0 0 | 0
So, for the given system, r ( A) r ( A | B) 3 number of unknowns. Hence the system has unique
solution. So the answer is 1.
Number of Solutions for homogeneous Linear System in terms of rank: Consider the
system of linear equations given by Eq. 1.5, where all b j 0 , 1 j m . In matrix form the system of
equation can be written as Ax O , where A a jk O is the coefficient matrix; x xk 1 n1 is
m n
the column vector. Now,
The system Ax O always has the trivial solution, i.e., x1 x2 xn 0 . Hence the system
is always consistent. [This point was asked in CS-1996 (1 mark)].
The system Ax O have non-trivial solution if and only if rank ( ) n . Hence, in this case, the
given system of equation has non-trivial solutions as well as trivial solution and hence the system
has infinitely many solutions. [This point was asked in CS-1996 (1 mark), CH-2009 (1 mark)].
Example 1.139 [EE-2008 (2 marks)]: A is a m n full rank matrix with m n and I is an identity
matrix. Let matrix A ( AT A) 1 AT . Then, which one of the following statement is FALSE?
(a) AA A A (b) ( AA ) 2 AA
(c) A A I
(d) AA A A
T
Solution (d): Given order of A is a m n order of A is a n m
For the order of A
A ( AT A) 1 AT (( n m
)( m n )) 1 (n m ) ( n n ) 1 (n m ) (n n
)(n m ) n m
equal equal
Option (a): order of AA A ( m n)(n m)( m n ) (m n ) order of A , so option (a) is correct.
equal equal
2 2 x1 0
Example 1.140 [EE-2013 (1 mark)]: The equation has
1 1 x2 0
(a) No solution (b) Only one solution x1 x2 0
(c) Non-zero solution (d) Multiple solutions
Solution (d): As we have homogeneous system of equation, whose coefficient matrix is:
2 2 2 2 1
A . So for r ( A) : A [By applying R2 R2 R1 ]. As we have only one
1 1 0 0 2
independent row so r ( A) 1 2 , so the given system of equation has non-trivial solutions as well as
trivial solution and hence the system has infinitely many solution.
Exercise: 1.4
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
T T
1. Express u 1 2 1 , as a linear combination of two vectors v1 1 2 3 and
T
v 2 1 0 2 ?
(a) u v1 2 v 2 (b) u 2 v1 v 2 (c) u v1 v 2 (d) Cannot be expressed
T T
2. Express u 1 2 0 , as a linear combination of two vectors v1 1 2 3 and
T
v 2 1 0 2 ?
(a) u v1 2 v 2 (b) u 2 v1 v 2 (c) u v1 v 2 (d) Cannot be expressed
T T T T
3. Let v1 1 2 , v 2 0 1 and v 3 1 1 and v 4 1 0 . Is v 4 can be expressed as a
linear combination of v1 , v 2 and v3 ?
(a) cannot be expressed (b) can be expressed in only one way
(c) can be expressed in infinite number of ways (d) None of these
T T
4. If v1 1 2 3 and v 2 1 0 2 then span {v1 , v 2 } is
(a) x y 2 z 0 (b) x 4 y 2 z 0 (c) 4 x y 2 z 0 (d) 4 x 2 y z 0
T T
5. i e1 1 0 and j e2 0 1 span k , where k _____.
T T
6. v1 1 1 and v2 2 1 span k , where k _____.
T T T k
7. The vectors v1 1 1 , v2 2 1 and v3 3 2 span , where k _____.
T
8. Which one of the following statement if correct for the given three vectors: v1 1 2 3 ,
T
v 2 1 0 2 .
(a) The vectors are mutually perpendicular (b) The vectors are linearly dependent
(c) The vectors are linearly independent (d) The vectors are unit vectors
T
9. Which one of the following statement if correct for the given three vectors: v1 1 1 0 ,
T T
v 2 1 0 1 and v 3 3 1 2 .
(a) The vectors are mutually perpendicular (b) The vectors are linearly dependent
(c) The vectors are linearly independent (d) The vectors are unit vectors
T
10. Which one of the following statement if correct for the given three vectors: v1 1 1 3 ,
T T T
v 2 1 3 1 , v 3 3 1 1 and v 4 3 3 3 .
(a) The vectors are mutually perpendicular (b) The vectors are linearly dependent
(c) The vectors are linearly independent (d) The vectors are unit vectors
3
11. The dimension of the plane x 2 z 0 in is _____.
T
12. Which one of the following is the basis for vector space V spanned by vectors v1 1 1 0 ,
T T T
v 2 0 1 1 , v 3 2 3 1 and v 4 1 1 1 ?
T T T T
(a) 2 1 1 0 (b) 2 1 1 0 (c) 2 1 1 0 (d) 2 1 1 0
a b 2c
2a 2b 4c d
13. The basis and dimension of the subspace W , where a, b, c, d are real are,
b c d
3a 3c d
_____ and _____ respectively.
1 1 0 1 1 2 1 1 0 1 1 2
2 2 1 2 2 4 2 2 1 2 2 4
(a) ,3 (b) ,3 (c) ,4 (d) ,4
0 1 1 0 1 1 0 1 1 0 1 1
3 0 1 3 0 3 3 0 1 3 0 3
T T T T
14. The pair of vectors: (i) 1 3 and 3 1 ; (ii) 1 3 and 10 10 3 ; (iii)
T T T T T
2 4 1 and 4 1 2 ; (iv) x y and y x ; (v) a b a b and
T
1 1 1 which are orthogonal to each other are:
(a) (i), (ii) and (iii) (b) (ii), (iii) and (iv) (c) (iii), (iv) and (v) (d) (ii), (iv) and (v)
T
15. The two vectors which are orthogonal to v 1 4 7 , and are not scalar multiples of each
other, are
T T T T
(a) 10 1 2 and 1 3 1 (b) 1 3 1 and 2 1 4
T T T T
(c) 3 1 1 and 10 1 2 (d) 3 1 1 and 2 1 4
T T
16. Angle (in degrees) between two vectors u 4 0 2 2 and v 2 0 1 1 is ____.
T
17. The orthogonal projection of y 1 4 3 onto the subspace spanned by vectors
T T
u 1 1 0 and u 1 1 0 is
T T T T
(a) 1 4 0 (b) 1 4 0 (c) 4 1 0 (d) 4 1 0
18. Suppose {v1 , v 2 , , v n } is a basis of V and {w1 , w 2 , , w m } is a linearly independent subset of
V , then which one of the following is correct?
(a) m n (b) m n (c) m n (d) m n
19. If {v1 , v 2 , , v m } is linearly dependent, then there exists _____ such that v j is a linear
combination of {v1 , v 2 , , v j1} .
(a) j m (b) j m (c) j m (d) j m
1 0 2 1
20. The rank of the matrix A 0 2 4 2 is _____.
0 2 2 1
3 6 1 1 7
21. The nullity of the matrix A 1 2 2 3 1 is _____.
2 4 5 8 4
1 2 1 0
22. The rank of matrix A 2 4 2 0 is _____.
4 8 4 0
23. Which of the following statement is not correct for a matrix A aij having rank 4?
7 6
(a) any 5 columns are dependent (b) some set of 4 columns are independent
(c) any 6 rows are dependent (d) any 4 rows are independent
1 5 1
24. The rank of matrix A 2 1 1 is _____.
3 6 2
25. The rank of a null matrix is always _____.
26. For a matrix A aij having six rows are dependent, if r be the rank of matrix A , then
6 7
which one of the following is true?
(a) 0 r 5 (b) 0 r 5 (c) 0 r 5 (d) 0 r 5
27. If a 5 7 matrix has rank 3, then which one of the following statement is not correct?
(a) every 5 5 sub-determinants is 0 (b) every 4 4 sub-determinant is zero
(c) every 3 3 sub-determinants is not 0 (d) at least one of all 3 3 sub-determinant is not 0
28. Which of the following statement is needed for determining the rank of matrix A aij ?
6 6
(a) One of its 3 3 sub-determinant is 6
(b) The 3 3 sub-determinant in north-east corner is 0 and in south-west corner is 8
(c) One of its 3 3 sub-determinants is 0, another of its 3 3 sub-determinants is 8 and all the
4 4 sub-determinants are 0
(d) One of its 3 3 sub-determinants is 0
a b
29. For the matrix A 2 , which one of the following statement is not correct?
0 a (b a )
(a) If a b , then rank ( A) 1 (b) If a b , then rank ( A) 1
(c) rank ( A) 2 for all values of a and b (d) If a b , then rank ( A) 2
30. What is the rank of matrix A aij , where aij i j 1 and n 1 ? _____.
n n
31. If A aij have rank r1 ; and AB abij have rank r2 , then which one of the following is
3 4 35
correct regarding rank ( r ) of matrix B ?
(a) r r1 r2 4 (b) r r1 r2 4 (c) r r2 r1 4 (d) r r2 r1 4
32. If matrices A aij , AB abij have full rank, then rank of B is?
3 4 35
1 1 1
34. What is the rank of the matrix A b c a c a b if any two of a, b, c are equal and is
bc ac ab
different from the third one? _____.
1 2 3 4
5 6 7 8
35. The rank of the matrix A is _____.
9 10 11 12
13 14 15 16
The eigenvalues of skew-Hermitian matrices are either purely imaginary or zero. The eigenvalues
of a skew symmetric matrices are either purely imaginary or zero; as skew-symmetric matrices
are skew-Hermitian matrices. [This point was asked in EC-2010].
Proof: Let A be a skew-Hermitian matrix, iA is Hermitian. Let be an eigenvalue of A
AX X , X 0 or (iA) X (i ) X i is an eigenvalue of iA . As iA is Hermitian and
its eigenvalues are real i is real, which is possible only if is zero or purely imaginary.
The eigenvalues of a orthogonal and unitary matrix are of unit modulus, i.e., 1 .
Proof: Let A be an orthogonal matrix AAT AT A I . Let be an eigenvalue of A
corresponding to an eigenvector X of matrix A , AX X , X 0 . Post-multiplying with
AT AT AX AT X IX AT X ( AT A I ) X AT X AT X X 1 is a
T T
eigenvalue of A . Since, the eigenvalues of matrix A and A are same
1 2 1 1 . Similar proof was also given for unitary matrices.
The eigenvalues of an idempotent matrix are either zero or unity.
Proof: Since A is an idempotent matrix A2 A . Let X be the eigenvector of the matrix A
corresponding to the eigenvalue AX X …(i), such that X O . On pre-multiplying (i)
by A A( AX ) A( X ) A2 X ( AX ) AX ( X ) X 2 X 2
( A A and
2
AX X ) ( ) X O 0 ( X O ) 0 or 1
2
23 13 23
Example 1.141 [CE-1999 (5 marks)]: Show that the matrix A 2 3 2 3 1 3 is orthogonal
1 3 2 3 2 3
and determine its Eigenvalues.
Solution: Any square matrix A of order n is said to be orthogonal, if AT A 1 AAT AT A I n .
2 3 1 3 2 3 2 3 2 3 1 3 1 0 0
AA 2 3 2 3 1 3 1 3 2 3
T
2 3 0 1 0 I 3 A is an orthogonal matrix.
1 3 2 3 2 3 2 3 1 3 2 3 0 0 1
(2 3) 13 23
Now for Eigenvalues, A I 0 2 3 (2 3) 13 0
13 23 (2 3)
5 11
9 3 6 2 6 9 0 (9 2 15 9)( 1) 0 1, i .
6 6
2 1 0 0
0 3 0 0
Example 1.142 [EC-2000 (1 mark)]: The Eigenvalues of the matrix are
0 0 2 0
0 0 1 4
(a) 2, –2, 1, –1 (b) 2, 3, –2, 4 (c) 2, 3, 1, 4 (d) None of these
Solution (b): For the given matrix if we apply the operation R4 R4 R3 2 , then we get upper
triangular matrix with diagonal elements 2, 3, –2, 4. So the Eigen values of the given matrix are
2, 3, 2, 4 ; so option (b) is correct.
[Similar question was also asked in CS-2002 (5 marks)]
2 1 0
Example 1.143 [CE-2000 5 marks]: A is a square matrix is given by A 1 2 1 . One of
0
1 2
the Eigenvalues of A is given to be equal to 2. Determine other Eigenvalues. Express the Eigenvector
T
corresponding to the lowest Eigenvalue in the form 1 a b and determine a and b .
Solution: Let the other two Eigenvalues are 2 and 3 . Since, sum of Eigenvalues of any non-
singular square matrix A trace of A sum of diagonal elements of A . Also, product of all
Eigenvalues of matrix A is equal to A . Hence, 2 2 3 2 2 2 2 3 4 …(i); and
223 2(4 1) 1(2 0) 4 23 2 …(ii). (i) and (ii) 2 (4 2 ) 2 22 42 2 0
2 2 2 3 2 2 . Thus minimum eigenvalue is min(2, 2 2, 2 2) 2 2 .
For Eigenvector corresponding to 2 2 , ( A I ) X 0
2 1 0 x 2 x1 x2 0 (i)
1
1 2 1 x2 0 x1 2 x2 x3 0 (ii) . (i) x2 2 x1 . Putting value of (i) in
x
0 1 2 3 x2 2 x3 0 (iii)
T
(iii) x3 x1 . Hence Eigenvector corresponding to 2 2 is X x1 2 x1 x1 .
T
Comparing this Eigenvector with 1 a b we have x1 1 a 2 x1 2 and b x1 1 .
[Similar questions were also asked in CE-2006, AE-2012 (2 marks)]
5 3
Example 1.144 [AG-2016 (1 mark)]: Eigen values of the matrix are
1 4
(a) –6.3 and –2.7 (b) –2.3 and –6.7 (c) 6.3 and 2.7 (d) 2.3 and 6.7
5 3
Solution (c): Let A and be the eigenvalues of A , then we must have A I 0 , i.e.
1 4
5 3 1 0 5 3 2
0 0 (5 )(4 ) 3 0 9 17 0 6.3, 2.7 .
1 4 0 1 1 4
[Similar questions were also asked in CE-2002, EE-1994, ME-2003, AG-2008, AG-2013, BT-
2013, MT-2012, AG-2013, EC-1998, PI-2016, MT-2016, BT-2016 (1 mark); CE-2001, CE-2004,
CE-2008, CE-2012, CS-2005, ME-1999, AE-2007, AG-2012, AG-2014, PI-2011 (2 marks)]
1 0 0 0
100 1 0 0
Example 1.145 [EE-2002 (1 mark)]: The determinant of the matrix is
100 200 1 0
100 200 300 1
(a) 100 (b) 200 (c) 1 (d) 300
Solution (c): Since the given matrix is lower triangular matrix, so its diagonal elements (1, 1, 1, 1) are
its Eigenvalues. So the determinant of the given matrix is the product of Eigenvalues which is 1.
1 1 3
Example 1.146 [ME-2004 (1 mark)]: The sum of the Eigen values of the matrix 1 5 1 is
3 1 1
(a) 5 (b) 7 (c) 9 (d) 18
Solution (b): Since, sum of eigenvalues of any square matrix A trace of A sum of diagonal
elements of A . Hence the sum of its eigenvalues is 1 5 1 7 .
[Similar questions were also asked in CE-2007, EE-1998, MN-2009, TF-2010, ME-2008, CE-
2014 (1 mark)]
Example 1.147 [CE-2005 (2 marks)]: Consider the system of equations Ann X n1 Bn1 . Let i , xi
be an Eigen-pair of an Eigenvalue and its corresponding Eigenvector for real matrix. Let I be a n n
unit matrix. Which one of the following statement is NOT correct?
(a) For a homogeneous n n system of linear equations ( A I ) X 0 having a non-trivial solution,
the rank of ( A I ) is less than n .
m
(b) For matrix A , m being a positive integer, (im , xim ) will be the Eigen-pair for all i
T 1 T
(c) If A A , then i 1 i (d) If A A , then i is real i
Solution (b): For a homogeneous n n system of linear equations ( A I ) X 0 having a non-
trivial solution we must have A I 0 the rank of ( A I ) is less than n . So option (a) is
correct. If is an eigenvalue of matrix A AX X . Now pre-multiplying by A
2 2 2 2
A X AX X Eigenvalue and Eigenvector pair of A is ( , X ) . Again pre-multiplying
3 3
by A A3 X 2 AX 3 X Eigenvalue and Eigenvector pair of A is ( , X ) . When pre-
multiplying the previous equation by A for m th we can easily get Am X m X Eigenvalue and
m m
Eigenvector pair of A is ( , X ) . Hence option (b) is not correct. For orthogonal matrices (
T T T 1
A A AA I A A ) the Eigenvalues are of unit modulus; so option (c) is correct. For
T
symmetric matrices ( A A ) the Eigenvalues are real; so option (d) is correct.
4 2
Example 1.148 [EC-2005 (2 marks)]: Given the matrix , the Eigenvector is
4 3
T T T T
(a) 3 2 (b) 4 3 (c) 2 1 (d) 1 2
4 2
Solution (c): For eigenvalues A I 0 0 (4 )(3 ) 8 0
4 3
2 20 0 ( 5)( 4) 0 5, 4 . For eigenvector of A corresponding to
4 ( 5) 2 x1 1 2 x1 x1 2 x2 0 (i)
5 0 0
4 3 ( 5) x2 4 8 x
2 4 x1 8 x2 0 (ii)
T
(i) and (ii) both x1 2 x2 . From the given options, choosing x2 1 x1 2 X 2 1 .
[Similar questions were also asked in ME-1994 (2.5 marks), CE-1998 (5 marks), EC-2007 (4
marks), EE-2010, ME-2010 (2 marks), EE-1995, TF-2007, ME-2014 (1 mark)]
3 2 2
Example 1.149 [EE-2005 (2 marks)]: For the matrix P 0 2 1 , one of the eigenvalues is
0 0 1
equal to –2. Which of the following is an eigenvector?
T T T T
(a) 3 2 1 (b) 3 2 1 (c) 1 2 3 (d) 2 5 0
Solution (d): For eigenvector ( X ) of P corresponding to eigenvalue 2 , we have
(P I ) X 0 .
3 ( 2) 2 2 x1 5 2 2 x1 5 x1 2 x2 2 x3 0 (i)
0 2 ( 2) 1 x 0 0
0 1 x3 0 x3 0 (ii)
2
0 0 1 ( 2) x3 0 0 3 x3 x3 0 (iii)
From (iii) we have x3 0 , from (1) we can choose x1 2 x2 5
[Similar question was also asked in XE-2012 (1 mark)]
4 2
Example 1.151 [EC-2006 (2 marks)]: For the matrix the Eigenvalue corresponding to the
2 4
T
Eigenvector 101 101 is:
(a) 2 (b) 4 (c) 6 (d) 8
Solution (c): The matrix and its Eigenvector is given and we have to find the corresponding
Eigenvalue. Since the Eigenvalue and its Eigenvector for a given matrix can be related as AX X .
4 2 101 101 606 101 101 101
6 6.
2 4 101 101 606 101 101 101
[Similar question was also asked in EE-1998, AE-2013 (1 mark), PI-2008 (2 marks)]
Example 1.152 [EC-2006 (2 marks)]: The Eigenvalues and the corresponding Eigenvectors of a
T T
2 2 matrix are given by 1 8 , v1 1 1 and 2 4 , v2 1 1 . The matrix is
6 2 4 6 2 4 4 8
(a) (b) (c) (d)
2 6 6 4 4 2 8 4
a b
Solution (a): Let the 2 2 matrix is A . Since the Eigenvalue and its Eigenvector for a
c d
given matrix can be related as AX X .
1 a b 1 1 a b 8 a b 8 (i)
For 1 8 , v1 , AX X 8
1 c d 1 1 c d 8 c d 8 (ii)
1 a 1 a b 4
b 1 a b 4 (iii)
For 2 4 , v2 , AX X 4
1 c d 1 1 c d 4 c d 4 (iv)
6 2
(i) and (iii) a 6 , b 2 . (ii) and (iv) c 2 , d 6 . Hence the matrix is .
2 6
3 2
Example 1.153 [ME-2006 (2 marks)]: Eigen value of a matrix A are 5 and 1. What are
2 3
2
the Eigen values of the matrix A AA ?
(a) 1 and 25 (b) 6 and 4 (c) 5 and 1 (d) 2 and 10
3 2
Solution (a): Eigenvalues of the matrix A are 5 and 1. We know that 1k , 2k , , nk are the
2 3
k 2 2
eigenvalues of A . So, eigenvalues of A are 52 and 1 , i.e., 25 and 1. Hence option (a) is correct.
[Similar question was also asked in CH-2008 (2 marks)]
1 2 1
Example 1.154 [AE-2007 (2 marks)]: The eigenvalues of the matrix A , where A are
0 3
(a) 1 and 1/2 (b) 1 and 1/3 (c) 2 and 3 (d) 1/2 and 1/3
2 1
Solution (d): Eigenvalue of the matrix A is determined as, A I 0
0 3
2 1
0 (2 )(3 ) 0 0 2, 3 . 1 1 ,1 2 , ,1 n are the eigenvalues of
0 3
A1 (provided A is non-singular). As A 6 1 5 0 The Eigenvalues of A1 are 1/2 and 1/3.
1 1 1
Example 1.155 [XE-2007 (1 mark)]: Let M 0 1 1 , then the maximum number of linearly
0 0 1
independent Eigenvectors of M is
(a) 0 (b) 1 (c) 2 (d) 3
Solution (d): As the given matrix is already in a row reduced echelon form and it has 3 non-zero
rows, so its rank is 3. Thus the number of linearly independent Eigenvectors of M is 3.
Example 1.158 [CS-2008 (2 marks)]: How many of the following matrices have an Eigenvalue 1?
1 0 0 1 1 1 1 0
0 0 , 0 0 , 1 1 and 1 1 .
(a) One (b) Two (c) Three (d) Four
1 0 1 0
Solution (a): Eigenvalues of , A I 0 0
0 (1 )( ) 0 1, 0 ;
0 0
0 1 1 1
Similarly, Eigenvalues of are 0, 0 ; Eigenvalues of 1 1 are 1 i ; Eigenvalues of
0 0
1 0 1 0
1 1 are 1 . Thus only the matrix 0 0 have Eigenvalue as 1.
p11 p12
Example 1.159 [EC-2008 (1 mark)]: All the four entries of the 2 2 matrix p are
p21 p22
non-zero, and one of its Eigenvalues is zero. Which of the following statements are true?
(a) p11 p22 p12 p21 1 (b) p11 p22 p12 p21 1
(c) p11 p22 p12 p21 0 (d) p11 p22 p12 p21 0
Solution (c): As, the product of eigenvalues of any square matrix gives determinant of that matrix. As
p11 p12
one of the eigenvalue is zero, so p 0 0 p11 p22 p12 p21 0
p21 p22
1 2
Example 1.160 [ME-2008 (2 marks)]: The Eigenvectors of the matrix are written in the
0 2
T T
form 1 a and 1 b . What is a b ?
(a) 0 (b) 1/2 (c) 1 (d) 2
Solution (b): As the given matrix is upper triangular so its Eigenvalues are 1, 2 . If X is the
Eigenvector corresponding to the Eigenvalue ( ) of the matrix A , then AX X .
1 2 x1 x1 x1 2 x2 x1 x1 2 x2 x1 (i)
For 1, AX X 1 2 x x (ii)
0 2 x2 x2 2 x2 x2 2 2
Example 1.166 [EE-2009 (1 mark)]: The trace and determinant of a 2 2 matrix are known to be –2
and –35 respectively. The Eigenvalues are
(a) –30 and –5 (b) –37 and –1 (c) –7 and 5 (d) 17.5 and –2
Solution (c): Let 1 and 2 are the two Eigenvalues of the given matrix, then 1 2 2 and
12 35 . From these two equations 1 ( 2 1 ) 35 12 21 35 0
1 5, 7 2 7,5
Example 1.167 [IN-2009 (2 mark)]: The Eigenvalues of a 2 2 matrix X are –2 and –3. The
1
Eigenvalues of matrix ( X I ) ( X 5I ) are
(a) –3 and –4 (b) –1 and –2 (c) –1 and –3 (d) –2 and –4
2 0 1 0
Solution (c): Let X diag( 2, 3) and I diag(1,1)
0 3 0 1
3 0
X 5I diag( 2, 3) diag(5,5) diag( 2 5, 3 5) diag(3, 2) ; and
0 2
1 0
X I diag( 2, 3) diag(1,1) diag(2 1, 3 1) diag( 1, 2)
0 2
For finding the inverse of ( X I ) , we have X I 2 and C11 2; C12 0; C21 0; C22 1
2 0 T 2 0 1 adj ( X I ) 1 0
Cij adj ( X I ) Cij (X I)
0 1 0 1 X I 0 1 2
1 1 0 3 0 3 0
( X I ) ( X 5I ) A (say)
0 1 2 0 2 0 1
3 0
So, for eigenvalue ( ) of A , we have A I 0 0 3, 1
0 1
2 3
Example 1.168 [CS-2010 (2 marks)]: Consider the following matrix: A . If the
x y
Eigenvalues of A are 4 and 8, then
(a) x 4, y 10 (b) x 5, y 8 (c) x 3, y 9 (d) x 4, y 10
Example 1.170 [XE-2010 (1 mark)]: Which one of the following matrices has the same Eigenvalues
1 2
as that of ?
4 3
3 4 1 4 4 2 2 4
(a) (b) (c) (d) 1 3
1 2 2 3 1 3
Solution (b): As the matrix in option (b) is the transpose of the given matrix, hence the matrix in
option (b) and the given matrix has same Eigenvalues.
2 a
Example 1.171 [AE-2011 (2 marks)]: Consider the matrix A where a and b are real
b 2
numbers. The two Eigenvalues of this matrix 1 and 2 are real and distinct 1 2 when
(a) a 0 and b 0 (b) a 0 and b 0 (c) a 0 and b 0 (d) a 0 and b 0
2 a
Solution (c): For Eigenvalues ( ) of the given matrix, we have A I 0 0
b 2
2 4 (4 ab) 0 . Now this equation has real and distinct values of if its discriminant is
2
0 (4) 4 1 (4 ab) 0 ab 0 a 0 & b 0 or a 0 & b 0 . So option (c) is correct.
4 0
Example 1.172 [CH-2012 (1 mark)]: Consider the following 2 2 matrix A . Which one
0 4
of the following vectors is NOT a valid Eigenvector of the above matrix?
T T T T
(a) 1 0 (b) 2 1 (c) 4 3 (d) 0 0
Solution (d): The given matrix is a diagonal matrix so its eigenvalues, 4, 4 . For eigenvector
4 0 x1 x1 4 x1 4 x1 4 x1 4 x1 (i)
corresponding to 4 , AX X 4 .
0 4 x2 x2 4 x2 4 x2 4 x2 4 x2 (ii)
(i) and (ii) x1 , x2 ; but both x1 and x2 cannot be zero at the same time, since any non-zero
vector, X , is said to be a eigenvector of matrix A , if there exist a number s.t. AX X ; where,
is an eigenvalue of matrix A corresponding to eigenvector X .
Example 1.173 [CS-2012 (1 mark)]: Let A be the 2 2 matrix with elements a11 a12 a21 1
19
and a22 1 . Then the Eigen values of the matrix A are
(a) 1024, –1024 (b) 1024 2 , 1024 2 (c) 4 2 , 4 2 (d) 512 2 , 512 2
1 1 1 1
Solution (d): A , for its Eigenvalues, A I 0 0
1 1 1 1
(1 )(1 ) 1 0 2 . Now, we know that if 1 , 2 ,, n are the eigenvalues of a
k
square matrix A , then 1k , 2k , , nk are the eigenvalues of A . So the Eigenvalues of
19 19 18
A ( 2) 2( 2) 512 2
3 5 2
Example 1.174 [EC-2013 (1 mark)]: The minimum Eigenvalue of the matrix A 5 12 7 is
2 7 5
(a) 0 (b) 1 (c) 2 (d) 3
3 5 2 3 3 2
By Applying
Solution (a): As the determinant of the given matrix is 5 12 7 5 5 7 0.
C 2 C2 C
3
2 7 5 2 2 5
Hence one of the eigenvalues of the matrix A is zero, as the product of eigenvalue of any matrix is
equal to the determinant of that matrix. Thus from the given options, the minimum eigenvalue is 0.
Example 1.175 [EE-2013 (2 marks)]: A matrix has Eigenvalues –1 and –2. The corresponding
T T
Eigenvectors are 1 1 and 1 2 respectively. The matrix is
1 1 1 2 1 0 0 1
(a) (b) (c) (d) 2 3
1 2 2 4 0 2
Solution (d): If X is the Eigenvector corresponding to the Eigenvalue ( ) of the matrix A , then
a11 a12
AX X . Let A .
a21 a22
1 a11 a12 1 1 a11 a12 1 a11 a12 1 (i)
For 1 , X , 1 1 a a 1 a a 1 (ii)
1 a21 a22 1 21 22 21 22
1 a a
11 12 1 1 11a 2a 12 2 a11 2a12 2 (iii)
For 2 , X , 2
2 a 2a 4
2 a21 a22 2 21 22 a21 2a22 4 (iv)
(i) and (iii) a11 0 and a12 1 ; similarly (ii) and (iv) a22 3 and a21 2 .
[Similar question was also asked in XE-2009 (2 marks)]
1 4
Example 1.176 [AG-2014 (1 mark)]: The eigenvalues of A are
2 3
(a) 2 i (b) 1, 2 (c) 1 2i (d) non-existent
Solution (c): For eigenvalues ( ) of the given matrix, we have
1 4 2
A I 0 0 (1 )( 3 ) 8 0 2 5 0 1 2i .
2 3
Example 1.177 [CS-2014 (1 mark)]: Which one of the following statements is TRUE about every
n n matrix with only real eigenvalues?
(a) If the trace of the matrix is positive and the determinant of the matrix is negative, at least one of
its eigenvalues is negative
(b) If the trace of the matrix is positive, all its eigenvalues are positive
(c) If the determinant of the matrix is positive, all its eigenvalues are positive
(d) If the product of the trace and determinant of the matrix is positive, all its eigenvalues are positive
Solution (a): Since, product of eigenvalues is equal to the determinant of any square matrix. If
determinant of the matrix is negative then we have at least one of the eigenvalue to be negative; and in
this case we may have trace of the matrix to be positive.
Example 1.179 [IN-2014 (2 marks)]: For the matrix A satisfying the equation
1 2 3 1 2 3
[ A] 7 8 9 4 5 6 , the eigenvalues are
4 5 6 7 8 9
(a) (1, j, j ) (b) (1,1, 0) (c) (1,1, 1) (d) (1, 0, 0)
Solution (c):
1 2 3 1 2 3 1 2 3 1 2 3 1 2 3 1 2 3
[ A] 7 8 9 4 5 6 A 7 8 9 4 5 6 A 7 8 9 7 8 9 A 1 .
4 5 6 7 8 9 4 5 6 7 8 9 4 5 6 4 5 6
(By applying R2 R3 in the right side of the determinant). So product of eigenvalues of A must be
–1. Hence, from the given options, we can say that option (c) is correct.
0 1 1
Example 1.180 [EE-2014 (2 marks)]: A system matrix is given as follows: A 6 11 6 .
6 11 5
The absolute value of the ratio of the maximum eigenvalue to the minimum eigenvalue is _____.
0 1 1
Solution: For Eigenvalues of the matrix A , we have A I 0 6 11 6 0
6 11 5
3 2
6 11 6 0 ( 1)( 2)( 3) 0 1, 2, 3 .
So maximum Eigenvalue is 1 and minimum Eigenvalue is 3 . Hence, the absolute value of the
ratio of the maximum eigenvalue to the minimum eigenvalue is 1 3 .
Example 1.181 [MA-2014 (1 mark)]: Let A M 3 () be such that det( A I ) 0 , where I denotes
the 3 3 identity matrix. If the tr ( A) 13 and det( A) 32 , then the sum of squares of the
eigenvalues of A is _____.
Solution: We have 3 3 , so let 1 , 2 , 3 are the eigenvalues of the matrix A which can be assumed
1 0 0
as: A 0 2 0 A I 0 (1 1)(2 1)(3 1) 0 .…(i) tr ( A) 13 1 2 3 13
0 0 3
…(ii); det( A) 32 12 3 32 …(iii). Now by putting values of (ii) and (iii) in (i), we get
Example 1.184 [CE-2016 (1 mark)]: If the entries in each column of a square matrix M add up to 1,
then an eigenvalue of M is
(a) 4 (b) 3 (c) 2 (d) 1
a b
Solution (d): Let we have 2 2 matrix M , such that a c 1 and b d 1 .
c d
a b
For eigenvalues, , of M, M I 0 0 ( a )( d ) bc 0
c d
2 (a d ) ad (1 d )(1 a) 0 2 ( a d ) a d 1 0 …(i). As (i) is satisfied only
by when 1 , so 1 is the required eigenvalue. Hence option (d) is correct.
In general: If the sum of the elements in each row or column of a square matrix is equal to k ,
then k is the eigenvalue of that matrix.
Example 1.185 [CS-2016 (1 mark)]: Two eigenvalues of a 3 3 real matrix P are (2 1) and 3.
The determinant of P is _____.
Solution: As one of the eigenvalue of real matrix P is (2 i) , so another will be (2 i ) ; and the
given third eigenvalue is 3. Thus the determinant of P will be (2 i )(2 i)(3) (4 1)(3) 15 .
Example 1.186 [CS-2016 (1 mark)]: Suppose that the eigenvalues of matrix A are 1, 2, 4. The
determinant of ( A1 )T is _____.
Solution: The determinant of matrix A is the product of all its eigenvalues.
So A (1)(2)(4) 8 . ( A1 )T A1 1 A 1 8 0.125 .
x
Example 1.187 [EC-2016 (1 mark)]: Consider a 2 2 square matrix A , where x is
unknown. If the eigenvalues of the matrix A are ( j ) and ( j ) , then x is equal to
(a) j (b) j (c) (d)
Solution: We know that the product of eigenvalues of a matrix gives the determinant of that matrix.
x 2 2 2
So we have, ( j )( j ) x x . Thus, option (d) is correct.
Example 1.189 [EE-2016 (1 mark)]: Consider a 3 3 matrix with every element being equal to 1. Its
only non-zero eigenvalue is _____.
Solution: For eigenvalues ‘ ’ of given square matrix A , we must have A I 0 . As
1 1 1 1 1 1
A 1 1 1 . So A I 0 1 1 1 0
1 1 1 1 1 1
3 2
(1 ){(1 ) 2 1} 1{(1 ) 1} 1{1 (1 )} 0 3 0 2 ( 3) 0 0, 3 .
Thus non-zero eigenvalue of matrix A is 3.
3
Example 1.190 [EE-2016 (1 mark)]: A 3 3 matrix P is such that, P P . Then the eigenvalues
of P are
(a) 1,1, 1 (b) 1, 0.5 j 0.866, 0.5 j 0.866 (c) 1, 0.5 j 0.866, 0.5 j 0.866 (d) 0,1, 1
Solution (d): If be an eigenvalue and x be the eigenvector of P , then Px x …(i)
On pre-multiplying both sides of (i), we get P 2 x Px ( x ) 2 x P 2 x 2 x …(ii)
Again pre-multiplying both sides of (ii), we get P 3 x 2 Px 2 ( x) 3 x P3 x 3 x
Px 3 x x 3 x ( 2 1) x 0 ( 2 1) 0 , as x 0 . Thus 0, 1 .
a b c
Example 1.191 [MA-2016 (2 mark)]: Let M b d e be a real matrix with eigenvalues 1, 0
c e f
and 3. If the eigenvectors corresponding to 1 and 0 are (1,1,1)T and (1, 1, 0)T respectively, then the
value of 3 f is equal to _____.
Solution: As the value of the dot product of the eigenvectors corresponding to any pair of different
eigenvalues of a n n real symmetric matrix A is always zero. So let ( x1 , x2 , x3 )T be the eigenvector
corresponding to the eigenvalue ‘3’; hence (1,1,1)T ( x1 , x2 , x3 )T 0 x1 x2 x3 0 …(i); also
(1, 1, 0)T ( x1 , x2 , x3 )T 0 x1 x2 0 …(ii); from (i) and (ii), we have x2 x1 , x3 2 x1 . Thus
eigenvector corresponding to eigenvalue of ‘3’ will be ( x1 , x2 , x3 )T ( x1 , x1 , 2 x1 )T x1 (1,1, 2)T or
(1,1, 2)T . Now we know that if X be the eigenvector corresponding to eigenvalue ‘ ’ of matrix M
, then we must have MX X . So
a b c 1 1 a b 0 a b (vi)
for 0 , X (1, 1, 0)T , b d e 1 0 1 b d 0 b d (vii)
c e f 0 0 c e 0 c e (viii)
a b c 1 1 a b c 1 2a c 1 (iii)
for 1 , X (1,1,1) , b d
T
e 1 1 1 b d e 1 2b c 1 (iv)
c e f 1 1 c e f 1 2c f 1 (v)
Example 1.192 [ME-2016 (1 mark)]: The condition for which the eigenvalues of the matrix
2 1
A are positive is
1 k
(a) k 1 2 (b) k 2 (c) k 0 (d) k 1 2
Solution (a): Let be the eigenvalues of A , then we must have A I 0 , i.e.
2 1 1 0 2 1
0 0 (2 )( k ) 1 0 (2 )( k ) 1 0
1 k 0 1 1 k
2 (k 2) (2k 1) 0 …(i); f ( ) 0 , where f ( ) 2 ( k 2) (2k 1)
We have to find values of k for which the roots of (i) are positive; so we must have
{(k 2)}2 4(1)(2k 1) 0 k 2 4k 8 0 ( k 2) 2 4 0 k …(i);
{( k 2) 2(1)} 0 k 2 …(iii); and (1) f (0) 0 (2k 1) 0 k 1 2 …(iv).
So from (ii), (iii) and (iv), we have k 1 2 . Hence option (a) is correct.
3
Solution: If 1 , 2 , 3 are eigenvalues of matrix A , then 13 , 23 , 33 are eigenvalues of matrix A ; and
2
312 , 322 , 332 are eigenvalues of matrix A . So (13 312 ), ( 23 322 ), (33 332 ) are the eigenvalues
of matrix ( A3 3A 2 ) . Hence {13 3(1) 2 },{( 1) 3 3( 1) 2 },{(3)3 3(3) 2 } 2, 4, 0 are the
3 2
eigenvalues of matrix ( A 3A ) .
As sum of eigenvalues of any matrix gives trace of that matrix, . trace( A 3 3A 2 ) 2 4 0 6 .
Example 1.195 [ME- Solution: For eigenvalues of the given matrix, we must have
2016 (2 marks)]: The 2 1 0
number of linearly
independent eigenvectors A I 0 0 2 0 0
of matrix 0 0 3
2 1 0 (2 ){(3 )(2 ) 0} 0 2, 2,3 . We know that if we have
A 0 2 0 is _____. n distinct eigenvalues then we have n linearly independent eigenvectors.
As we have only two eigenvalues, i.e. 2, 3 , so we have only two
0 0 3 linearly eigenvectors. So, answer is two.
5 3
Example 1.196 [PE-2016 (2 marks)]: Consider the matrix, M . The normalised eigen-
3 5
vector corresponding to the smallest eigen-value of the matrix M is
T T T T
(a) 3 2 1 2 (b) 3 2 1 2 (c) 1 2 1 2 (d) 1 2 1 2
5 3
Solution (c): For eigenvalues ( ) of M , we must have M I 0 0
3 5
(5 )(5 ) (3)(3) 0 2 10 16 0 ( 2)( 8) 0 2,8
So smallest eigenvalue of the matrix M is 2 .
Hence eigenvector ( X ) corresponding to eigenvalue 2 is given by MX X
5 3 x1 x1 5 x1 3 x2 2 x1 5 x1 3x2 2 x1 x2 x1 ...(i)
3 5 x2 x2 3 x1 5 x2 2 x2 3 x1 5 x2 2 x2 x2 x1 ...(ii)
T
Thus X x1 x1 is the required eigenvector corresponding to eigenvalue 2 .
Hence from given options, option (c) is the required normalised eigenvector.
[Similar question was also asked in ME-2012, PI-2012 (2 marks)]
3 4
Example 1.197 [TF-2016 (2 marks)]: The eigenvalues and eigenvectors of are
4 3
T T T T
(a) 5 and 1 2 , 2 1 respectively (b) 3 and 1 2 , 2 1 respectively
T T T T
(c) 4 and 1 2 , 2 1 respectively (d) 5 and 1 1 , 2 1 respectively
3 4 3 4
Solution (a): For eigenvalues ‘ ’ of A , A I 0 0
4 3 4 3
(3 )( 3 ) 16 0 2 25 5 .Eigenvector corresponding to 5 is given by
3 4 x1 x1 3 x1 4 x2 5 x1 x1 2 x2 x1 2 x2 2
Ax x : 5 x x2
4 3 x2 x2 4 x1 3x2 5 x2 x1 2 x2 x2 x2 1
Eigenvector corresponding to 5 is given by:
3 4 x1 x1 3 x1 4 x2 5 x1 x2 2 x1 x1 x1 1
5 x x1
4 3 x2 x2 4 x1 3 x2 5 x2 x2 2 x1 x2 2 x1 2
[Similar questions were also asked in CS-1993, EC-1993 (1 mark)]
Block matrix: A block matrix is a matrix that is defined using smaller matrices called blocks. For
0 2 3 9
0 2 2 0 1 3
3 9 5 0
e.g., if A 2 0 , B
1 3 , C 7 4 , D 0 5 then M 4 6 5 0 then matrix
4 6
7 4 0 5
A B
M , where A, B , C , D are themselves matrices, is a block matrix.
C D
A B
If we have a block matrix defined by M , where A, B , C , D are of same order, and
C D
also, if C and D commute (i.e., CD DC ); then det M det( AD BC )
Example 1.198 [CS-2007 (2 marks)]: Let A be a 4 4 matrix with Eigenvalues –5, –2, 1, 4. Which
A I
of the following is an Eigenvalue of M , where I is the 4 4 identity matrix?
I A
(a) –5 (b) –7 (c) 2 (d) 1
A I I
Solution (c): Let be the Eigenvalue of the matrix M then M I 0 0.
I A I
As we have block matrix and it is given that A and I both are 4 4 matrices, so ( A I ) is also
A I I 2 2
4 4 matrix; also I ( A I ) ( A I ) I ; hence 0 ( A I ) I 0
I A I
( A I I )( A I I ) 0 A I I A I I 0 A I I 0 or A I I 0 .
As, A be a 4 4 matrix with Eigenvalues –5, –2, 1, 4 so matrix A can be written as
5 0 0 0
0 2 0 0
A diag ( A) ( 5, 2,1, 4) ; since I diag (1,1,1,1) and I diag ( , , , ) .
0 0 1 0
0 0 0 4
So, A I I diag (5 1, 2 1,1 1, 4 1)
A I I 0 ( 4 )(1 )(2 )(5 ) 0 4, 1, 2,5
Also, A I I diag ( 5 1, 2 1,1 1, 4 1)
A I I 0 ( 6 )( 3 )( )(3 ) 0 6, 3, 0, 3
Proof: All we have to show is that x1 , x2 ,, x n are linearly independent. Suppose they are not. Let r
be the largest integer such that {x1 , x2 ,, xr } is a linearly independent set. Then r n and the set
{x1 , x2 ,, xr , x r1} is linearly dependent. Thus there exists scalars c1 , c2 ,, cr 1 , not all zero, s.t.
c1x1 , c2 x2 ,, cr 1xr 1 0 (1.31)
Multiplying both sides of Eq. 1.14 by A and using A x j j x j , we obtain,
c11 x1 c22 x2 cr 1r 1 xr 1 0 (1.32)
To get rid of the last term, we subtract r 1 times Eq. 1.14 from 1.15, obtaining
c1 (1 r 1 )x1 c2 (2 r 1 )x2 cr (r r 1 )xr 0 (1.33)
Example 1.199 [XE-2009 (1 mark)]: Let A and B be two similar square matrices of order two. If 1
and –2 are the Eigenvalues of A , then the trace of B is
(a) –2 (b) –1 (c) 1 (d) 2
Solution (b): As A and B are two similar square matrices Eigenvalues of A and B are same.
So tr ( B) 1 2 1 . Hence option (b) is correct.
Diagonalization of Matrices: A diagonal matrix is a n n matrix with non-zero entries only along
the main diagonal. Diagonal matrices are particularly convenient for eigenvalue problems since the
eigenvalues of a diagonal matrix coincide with the diagonal entries aii and the eigenvector
corresponding to the eigenvalue aii is just the i th coordinate vector. If any square matrix is not
diagonal matrix then diagonalisation is the process of finding a corresponding diagonal matrix. Once
we have a diagonal matrix, we can easily find the eigenvalues or eigenvectors.
An n n matrix A is diagonalizable if and only if it has n linearly independent eigenvectors.
Proof: Suppose A has n linearly independent eigenvectors. Then the matrix C formed by using
these eigenvectors as column vectors will be invertible (since the rank of C will be equal to n ).
On the other hand, if A is diagonalizable then, by definition, there must be an invertible matrix
1
C such that D C AC is diagonal. But the columns vectors of C must coincide with the
eigenvectors of A . Since C is invertible, these n column vectors must be linearly independent.
Hence, A has n linearly independent eigenvectors.
If a n n matrix A is diagonalizable then it does not mean that matrix A has distinct
Eigenvalues; as we have proved earlier that an Eigenvalue of a matrix can correspond to different
Eigenvectors. However, if n n matrix A has n distinct Eigenvalues then A is diagonalizable.
T
If A is a real symmetric matrix (i.e., A A ) then (i) it has only real Eigenvalues (ii) it is always
diagonalizable (iii) it has orthogonal eigenvectors [This point was asked in AE-2014 (1 mark)].
Let A be a n n matrix, let 1 , 2 ,, n be a set of n scalars, and let v1 , v 2 ,, v n be a set of n
th
vectors. Let C be n n matrix formed by using v j for j column vector, and let D be the
n n matrix whose diagonal entries are 1 , 2 ,, n . Then AC CD if and only if 1 , 2 ,, n
are eigenvalues of A and each v j is an eigenvector of A corresponding to the eigenvalue j .
| | | |
Proof: Under the hypothesis AC A v1 v n A v1 A v n
| | | |
| | 1 0 | | A v1 1v1
CD v1 v n 1 v1 n v n and so, AC CD
| | 0 n | | A vn n v n
and vice-versa. Now suppose AC CD , and the matrix C is invertible. Then we can write
1
D C AC and so we can think of the matrix C as converting A into diagonal matrix. Hence,
An n n matrix A is diagonalizable if there is an invertible n n matrix C such that C 1 AC is
a diagonal matrix. The matrix C is said to diagonalize A .
1 b a
Example 1.200 [XE-2007 (2 marks)]: Let M 0 2 c , where a , b and c are real numbers.
0 0 1
Then M is diagonalizable
(a) for all values of a , b and c (b) only when bc a
(c) only when b c a (d) only when bc a
Solution (d): We have to find the condition for which the given matrix has 3 linearly independent
T
eigenvectors. Let an eigenvector of M is given as X x1 x2 x3 . As the Eigenvalues of the
given matrix are 1, 2,1 . For Eigenvector corresponding to 2 , we have MX X
1 b a x1 x1 x1 bx2 ax3 2 x1 (i) (iii) x3 0
0 2 c x2 2 x2 2 x2 cx3 2 x2 (ii) (ii) x2 k (say) .
0 0 1 x3 x3 x3 2 x3 (iii) (i) x1 bk
T T
So X 1 bk k 0 k b 1 0 (which is one independent Eigenvector).
1 b a x1 x1
For Eigenvector corresponding to 1, we have MX X 0 2 c x2 1 x2
0 0 1 x3 x3
x1 bx2 ax3 x1 (iv) ; 2 x2 cx3 x2 (v) ; x3 x3 (vi) . Let x3 q (say); then from (ii), we get
x2 cq ; and putting values of x3 and x2 in (i), we get ( a bc) q 0 (vii) ; so we can let x1 p
T T
(say); hence X 2 p cq q . Now from (vii) if ( a bc ) 0 q 0 X 2 p 0 0
we have only one independent Eigenvector vector and so we have total two independent
Eigenvectors; hence matrix M is not diagonalizable. On the other hand, if from (vii) if
T
q 0 ( a bc) 0 X 2 p cq q we can have two independent Eigenvector vector
and so we have total three independent Eigenvectors; hence matrix M is diagonalizable. So for
matrix M to be diagonalizable we have the condition ( a bc ) 0 .
Example 1.201 [XE-2011 (1 mark)]: For any positive number a and b , the matrix
T
P 1 a b 4 5 6 is
(a) Orthogonal (b) diagonalizable (c) non-singular (d) of rank 2
4 5 6 4 5 6 1 1 1
Solution (b): P 4a 5a 6a P 4a 5a 6a 4 5 6 a a a (By taking common 4,
4b 5b 6b 4b 5b 6b b b b
5, 6 from C1 , C2 , C3 , respectively); so P 0 , as the two columns are identical. So P is a singular
By Applying 4 5 6
matrix. R2 R2 aR1 , P 0 0 0 P has only one independent row so rank of P is 1.
R R bR
3 3 1
0 0 0
4 5 6 4 4a 4b 77 77 a 77b
Also, PPT 4a 5a 6a 5 5a 5b 77 a 77 a 2 77 ab I ; P is not orthogonal
4b 5b 6b 6 6a 6b 77b 77 ab 77b 2
matrix. As options (a), (c) and (d) are incorrect; so we have only one option (b) which is correct.
1 0 1
Example 1.202 [XE-2013 (2 marks)]: For the matrix M 0 1 1 . Consider the following
1 1 2
1
statements: (A) The characteristic equation of M is 3 0 . (B) M does not exist. (C) The
matrix M is diagonalisable. Which of the above statement are true?
(a) A, B and C (b) A and C but not B (c) A and B but not C (d) B and C but not A
1 0 1
Solution (a): For characteristic equation, we have M I 0 0 1 1 0
1 1 2
3
(1 ){(1 )( 2 ) 1} 1{0 (1 )} 0 0 . So statement (A) is correct. As the
characteristic equation is 3 0 ( 2 1) 0 0, 1 one of the Eigenvalue is zero so
determinant of matrix M is zero, hence M is a singular matrix so it is not invertible; so statement
T
(B) is correct. Now for finding the Eigenvector, we have MX X , where X x1 x2 x3 .
1 0 1 x1 x1 x1 x3 0 (i)
Now, Eigenvector corresponding to 0 0 1 1 x2 0 x2 x2 x3 0 (ii)
1 1 2 x3 x3 x1 x2 2 x3 0 (iii)
T T
(i) and (ii) x1 x2 x3 k (say) X 1 k k k k 1 1 1
1 0 1 x1 x1 x1 x3 x1 (i)
Eigenvector corresponding to 1 0 1 1 x2 1 x2 x2 x3 x2 (ii)
1 1 2 x3 x3 x1 x2 2 x3 x3 (iii)
T T
(i) and (ii) x3 0 and (iii) x2 x1 p (say) X 2 p p 0 p 1 1 0
1 0 1 x1 x1 x1 x3 x1 (i)
Eigenvector corresponding to 1, 0 1 1 x2 1 x2 x2 x3 x2 (ii) . So
1 1 2 x3 x3 x1 x2 2 x3 x3 (iii)
(i) x1 x3 2 ; (ii) x2 x3 2 ; (iii) x3 x3 r (say)
T T
X 3 r 2 r 2 r r 1 2 1 2 1 . So we have three linearly independent Eigenvectors of
matrix M so it is a diagonalizable matrix. So statement (C) is also correct.
1 1
Example 1.203 [XE-2016 (2 marks)]: Let M . Which of the following is correct?
0 1
(a) Rank of M is 1 and M is not diagonalizable
(b) Rank of M is 2 and M is diagonalizable
(c) 1 is the only eigenvalues and M is not diagonalizable
(d) 1 is the only eigenvalue and M is diagonalizable
1 1
Solution (c): For eigenvalues ‘ ’ of M , M I 0 0 (1 )(1 ) 0 0
0 1
1,1 . As eigenvalues of M are not distinct, so M is not diagonalizable. So option (c) is
correct.
x1 x1 1
can choose X 2 x1 . So matrix P can be chosen as the independent
x2 2 x1 2
1 1 1
Eigenvectors of matrix A , i.e., P . Now we have to find P , as P 2 1 1 and
1 2
T T
1
Cij
adj ( P ) 1 2 1 2 1
C11 2; C12 1; C21 1; C22 1 . So P
P P P 1 1 1 1
1 1 e 1 0 2 1 2e 1 e 2 e 1 e 2
Now, e A Pe D P 1 e A 2
1 2
1 2 0 e 1 1 2e 2e e 1 2e 2
Cayley-Hamilton Theorem: In linear algebra Cayley-Hamilton theorem states that every square
matrix satisfies its own characteristic equation.
In other words if A be any n n matrix and I n be n n identity matrix, then the characteristic
polynomial of A is defined as p ( ) A I . Cayley-Hamilton Theorem states that substituting
the matrix A for in this polynomial result in zero or null matrix.
The powers of A , obtained by substitution from powers of , are defined by repeated matrix
0
multiplication; the constant term of p ( ) gives a multiple of the power A , which power is
defined as the identity matrix. The theorem allows
n
A to be expressed as a linear combination of the lower matrix powers of A .
any matrix polynomial in A of order n n as a polynomial of degree n 1 in A .
inverse of any square matrix as a linear combination of the matrix A .
Statement for Linked Answer Question 1.207 and 1.208: Cayley-Hamilton Theorem states that a
3 2
square matrix satisfies its own characteristic equation. Consider a matrix A .
1 0
Example 1.207 [EE-2007 (2 marks)]: A satisfies the relation
(a) A 3I 2 A1 0 (b) A2 2 A 2 I 0 (c) ( A I )( A 2 I ) 0 (d) exp( A) 0
3 2 2
Solution (a): A I 0 0 3 2 0 . Then by Cayley-Hamilton
1 0
Theorem, substituting A for , A2 3 A 2 I 0 . Now multiplying both sides of the above
1
relation with A , A1 A2 3 A1 A 2 A1 I 0 A 3I 2 A 1 0 .
9
Example 1.208 [EE-2007 (2 marks)]: A equals
(a) 511A 510 I (b) 309 A 104 I (c) 154 A 155I (d) exp(9 A)
2 2
Solution (a): A 3 A 2 I 0 A 2 I 3 A .
A4 A2 A2 (2I 3 A)(2 I 3 A) 4I 12 A 9 A2 4I 12 A 9(2I 3 A) 15 A 14I
A8 A4 A4 (15 A 14I )(15 A 14I ) 225 A2 420 A 196I
A8 225(2I 3 A) 420 A 196I 255 A 254I
A9 A8 A 255 A2 254 A 255(2 I 3 A) 254 A 511A 510I A9 511A 510 I .
[Similar questions were also asked in EC-2012, EE-2012, IN-2012 (2 marks)]
Example 1.209 [EE-2007 (2 marks)]: The linear operation L( x) is defined by the cross product
T T
L( x ) b x , where b 0 1 0 and x x1 x2 x3 are three dimensional vectors. The 3 3
T
matrix M of this operation satisfies L( x ) M x1 x2 x3 . Then the Eigenvalues of M are
(a) 0, 1 , –1 (b) 1, –1, 1 (c) i , i , 1 (d) i , i , 0
iˆ ˆj kˆ x3
Solution (c): L ( x) b x 0 1 0 x3iˆ 0 ˆj x1kˆ . In matrix notation, L( x) 0 .
x1 x2 x3 x1
x1 x3 m11 m12 m13 x1 m11 m12 0; m13 1
0 0 1
L ( x ) M x2 0 m21 m22 m23 x2 m21 m22 m23 0 M 0 0 0 .
x3 x1 m31 m32 m33 x3 m31 1; m32 m33 0 1 0 0
0 0 1
For Eigenvalues of M , M I 0 0 0 0 0 3 0 1, i .
1 0 0
1
Example 1.210 [XE-2007 (2 marks)]: Let M be a 2 2 matrix with eigenvalue 1 & 2. Then M is
(a) ( M 3I ) 2 (b) (3I M ) 2 (c) (3I M ) 2 (d) ( M 3 I ) 2
Solution (c): As Eigenvalues of 2 2 matrix M are 1 and 2, so its characteristic equation can be
2
given as ( 1)( 2) 0 3 2 0 . Replacing M by , we get characteristic polynomials,
1
as, M 2 3M 2 I 0 . Now multiplying both of the characteristic polynomial with M ,
1 2 1 1 1
M M M 3M M 2 I 0 M 3 I 2 M 0
1 2 1
( M M M MM IM M ) M 1 (3I M ) 2 .
2 2 2 2
(a) ( P P 2I ) (b) ( P P I ) (c) ( P P I ) (d) ( P P 2I )
Solution (d): By Cayley-Hamilton Theorem, substituting P for , P3 P 2 2 P I 0 .
Pre-multiplying both sides by P P 1 P 3 P 1 P 2 2 P 1 P P 1 I 0
1
P 2 P 2I P 1 0 P 1 ( P 2 P 2 I )
1 1
Example 1.212 [XE-2010 (1 mark)]: Find P8 2 P 7 2 P 6 4 P5 3P 4 6 P 3 2 P 2 , if P ,
1 1
(a) P (b) 2 P (c) 3 P (d) 4 P
Solution (d): The characteristic equation for the given matrix can be found by P I 0
1 1 2 2
0 (1 ) 1 0 2 0 . Using Cayley-Hamilton Theorem, substituting P
1 1
2 2 3 2 2
for , P 2P 0 P 2P . Now P P P 2 P P 2P 4P ,
4 2 2 2 5 6 7 8
P P P 2 P 2 P 4 P 8P , and similarly, P 16 P , P 32 P , P 64 P , P 128 P .
So, P8 2 P 7 2 P 6 4 P 5 3P 4 6 P 3 2 P 2 128P 128 P 64 P 64 P 24 P 24 P 4 P 4 P
1
Example 1.213 [XE-2011 (1 mark)]: The Eigen values of a 3 3 matrix P are 2, 2, –1. The P is
equal to
(a) (3P P 2 ) 4 (b) ( P 2 2 P ) 2 (c) ( P 2 3P) 2 (d) ( P 2 2 P ) 4
Solution (a): The characteristic equation in is given as: ( 2)( 2)( 1) 0
( 2 4 4)( 1) 0 3 3 2 4 0 . Hence by Cayley-Hamilton Theorem, substituting P
1
for , P 3 3P 2 4 I 0 . Pre-multiplying by P , we get P 1 P 3 3P 1 P 2 4 P 1 I 0
1
2
P 3P 4 P
1
0 P (3P P 2 ) 4 .
3 3
Example 1.214 [AE-2014 (2 marks)]: If [ A] . Then [ A]2 7[ A] 3[ I ] is
3 4
(a) 0 (b) –324 (c) 324 (d) 6
Solution (a): The characteristic equation of the given matrix can be found as: A I 0
3 3 2
0 (3 )(4 ) 9 0 7 3 0 . Using Cayley-Hamilton Theorem,
3 4
substituting A for , A2 7 A 3I O ( A2 7 A 3I ) O O 0 .
2 0 3
Example 1.215 [MA-2014 (2 marks)]: Let A 3 1 3 . A matrix P such that P AP is a
1
0 0 1
diagonal matrix, is
1 1 1 1 1 1 1 1 1 1 1 1
(a) 0 1 1 (b) 0 1 1 (c) 0 1 1 (d) 0 1 1
1 1 0 1 1 0 1 1 0 1 1 0
Solution (a): For Eigenvalues of the given matrix, we have A I 0
2 0 3
T
3 1 3 0 (2 )( 1 ) 2 1, 2 . Now let X x1 x2 x3 is a
0 0 1
Eigenvector corresponding to Eigenvalue of matrix A , then AX X .
a11 a12
Solution (d): Let R
a21 a22
1 a11 a12 1 1 a11 1 a11 1
For 1 1 and V1 , RV1 1V1 1
0 a21 a22 0 0 a21 0 a21 0
1 a11 a12 1 1 a11 a12 3 a12 4
For 2 3 and V2 , RV2 2 V2 3
1 a21 a22 1 1 a21 a22 3 a22 3
1 4 3 1 0
So, R . If we apply R1 R1 R2 on matrix R , we get R .
0 3 4 0 3
1
As R and R are similar matrices because R P RP . So option (d) is correct.
1 4 2 1
8. 2 is an eigenvector of matrix A 2 0 1 . What is the corresponding eigenvalue? _____.
2 2 2 3
4 7 1
9. What is the sum of eigenvalues of the matrix A 0 3 8 ?
0 0 2
2 5 6
10. Given that 1 is an eigenvalue of the matrix A 1 0 0 . The other two eigenvalues are
0 1 0
(a) 2 and 3 (b) 2 and 3 (c) 0 and 3 (d) 2 and 0
4 1 6
11. The characteristic equation of the matrix A 2 1 5 is
2 1 0
(a) 3 5 2 4 0 (b) 3 5 2 21 42 0
(c) 3 5 2 4 0 (d) 3 5 2 14 0
1 2 4
12. Given that 5 is an eigenvalue of 1 4 8 , which of the following system of equations
0 1 1
should be solved to find the corresponding eigenvectors?
4 3 9 x 0 4 2 4 x x
(a) 6 1 3 y 0 (b) 1 1 8 y y
5 4 6 z 0 0 1 6 z z
4 2 4 x 0 1 2 4 x 5
(c) 1 1 8 y 0 (d) 1 4 8 y 5
0 1 6 z 0 0 1 1 z 5
13. A particular 3 3 matrix A has an eigenvalue of 1 . The matrix A I T
(a) 2t 0 t
1 0 2 T
3 2 12 1 1
16. If P , A and Q PAPT , then PT Q 2005 P
1 2 3 2 0 1
1 2005 1 2005 1 0 1 0
(a) (b) (c) (d)
0 1 2005 1 2005 1 0 1
1 0 0
17. If A 0 1 1 and 6A1 A2 cA dI , then ( d c ) _____
0 2 4
2 n
18. Matrix A is such that A 2 A I , where I is the identity matrix, then for n 2 , A
(a) 2n 1 A ( n 1) I n 1
(b) 2 A I (c) nA (n 1) I (d) nA I
0 a c
19. Let A and ( A I )50 50 A , then ( a b c d ) _____.
0 0 b d
i i 1 1
20. If A , where i 1 , and B , then A8 kB , where k _____.
i i 1 1
21. If P is an orthogonal matrix and Q PAPT and x PT Q1000 P , then x 1 is, where A is
involuntary matrix.
1000 (d) None of these
(a) A (b) I (c) A
22. If A and B are square matrices of same order and A is non-singular, then for a positive integer
1 n
n , ( A BA)
n n n n n n 1 n
(a) A B A (b) A B A (c) A B A (d) n( A1 BA)
1
23. If A and B are two square matrices such that B A BA , then ( A B ) 2
2 2 (b) O 2 2
(a) A B (c) A 2 AB B (d) A B
24. If one of the eigenvalues of a square matrix A is zero, then
(a) det A must be non-zero (b) det A must be zero
(c) adjA must be a zero matrix (d) tr ( A) must be zero
0 1 8 6 4 2 T
25. If A and ( A A A A I ) B 0 11 , where I is a 2 2 identity matrix, then
3 0
the product of all elements of matrix B is _____.
26. If the matrix A has as an eigenvalue with corresponding eigenvector x . The non-singular
matrix Q is of the same order as A . Then
(a) Q x is an eigenvector of the matrix B , where B QAQ 1 , and 1 is the corresponding
eigenvalue
(b) Q x is an eigenvector of the matrix B , where B QAQ 1 , and is the corresponding
eigenvalue
(c) B x is an eigenvector of the matrix Q , where B QAQ 1 , and is the corresponding
eigenvalue
(d) B x is an eigenvector of the matrix Q , where B QAQ 1 , and 1 is the corresponding
eigenvalue
27. If X is a diagonalizable n n matrix with only one eigenvalue , then
(a) X I n (b) I n X (c) I n n X (d) X n In
2 4 6 2 0 0
28. If matrices A 0 2 2 and B 2 6 0 , then
0 0 4 3 2 1
Answer Keys
Answer Keys: Exercise: 1.1
1 2 3 4 5 6 7 8 9 10 11 12 13 14 15
d 3 a b b 0 1 c 4 9 0 0 2 0 0
16 17 18 19 20 21 22 23 24 25 26 27 28 29 30
4 0 2 1 8 1 b b d d d 2.4 5 0 c
31 32 33 34 35 36 37 38 39 40 41 42 43 44 45
b 13 2 0 9 d d 36 1 c 256 d 4 c c
2
[AE – 2017 (1 mark)]: The value of the integral I {( x y ) dx x dy} , with C the boundary of
C
the square 0 x 2 ; 0 y 2 is _____.
Solution: Here F1 ( x, y ) x y and F2 ( x, y ) x 2 , so that F1 y 1 and F2 x 2 x . Thus the
line integral transforms into an easy double integral by using Green’s theorem in a plane as:
I {( x y ) dx x 2 dy} {(F2 x ) (F1 y )}dx dy , where R is the region around a square
C R
2 2
0 x 2 ; 0 y 2 . Thus I {2 x 1}dx dy dy (2 x 1) dx 2 (2 2 2) 12 .
R 0 0
[AE – 2017 (1 mark)]: Let v (t ) be a unit vector that is a function of the parameter t . Then
v {d v dt} _____.
Solution: Let v (t ) x (t ) i y (t ) j z (t ) k , as v (t ) 1 {x (t )}2 { y (t )}2 { z (t )}2 1
2 2 2 d{v (t )} d{ x(t )} d { y (t )} d {z (t )}
{x (t )} { y (t )} {z (t )} 1 . Now i j k; so
dt dt dt dt
d {x (t )} d{ y (t )} d{ z (t )} d{ x(t )}2 d { y (t )}2 d {z (t )}2
v {d v dt} x(t ) y (t ) z (t )
dt dt dt 2dt 2dt 2dt
1 d 1 d
v {d v dt } [{x (t )}2 { y (t )}2 { z (t )}2 ] [1] 0 .
2 dt 2 dt
[BT – 2017 (1 mark)]: The surface area (in m2) of the largest sphere that can fit into a hollow cube
with edges of length 1 meter is _____. Give data: 3.14 .
Solution: The diameter of the largest sphere will be equal to the edge of length 1 meter; so radius of
the sphere is 0.5 meter. Hence the largest surface area of the sphere is
A 4 r 2 4 3.14 (0.5) 2 3.14 m2.
2 2
[AG – 2017 (2 marks)]: I ( a x ) dx is
Solution: At x 0 , the given limit takes the form of 0 0 , so using L’Hospital rule, we have
lim{(sin x ) x} lim{(cos x ) 1} 1 .
x 0 x 0
[CE – 2017 (1 mark)]: Let x be a continuous variable defined over the interval ( , ) and
x
f ( x ) e x e . The integral g ( x) f ( x ) dx is equal to
x x x
(a) e e (b) e e (c) e e (d) e x
x x
Solution (b): g ( x) e x e dx e x e e dx ; now let e
x x
t e dx dt , so
e x
g ( x) et dt et c e c.
tan x
[CE – 2017 (1 mark)]: lim 2 is equal to _____.
x 0 x x
Solution: At x 0 , the given limit takes the form of 0 0 , so using L’Hospital rule, we have
tan x sec 2 x
lim lim
x0 1
x 0 x 2 x
2x 1
[CE – 2017 (1 mark)]: Let w f ( x, y ) , where x and y are functions of t . Then, according to the
chain rule, dw dt is equal to
dw dx dw dt w x w y w dx w dy dw x dw y
(a) (b) (c) (d)
dx dt dy dt y t x t x dt y dt dx t dy t
dw w dx w dy
Solution (c): Using chain rule, we have
dt x dt y dt
[CE – 2017 (2 marks)]: The tangent to the curve represented by y x ln x is required to have 45o
inclination with the x axis. The coordinates of the tangent point would be
(a) (1, 0) (b) (0,1) (c) (1,1) (d) ( 2, 2)
Solution (a): dy dx 1 ln x tan 45o 1 ln x 1 ln x 0 x 1 y 1(ln1) 0 .
1
1 (sin x)2
[CE – 2017 (2 marks)]: Consider the following definite integral: I dx . The value of
0 2
1 x
the integral is
(a) 3 24 (b) 3 12 (c) 3 48 (d) 3 64
1 1 1
Solution: Let sin 1 x t dx dt ; also at x 0, t sin 0 0 and at x 1, t sin 1 ;
2 2
1 x
1 2
1 (sin x)2 2 2 t3 3
so I dx t dt .
0 2 0 3 24
1 x 0
Solution: At x 0 , the given limit takes the form of 0 0 , so using L’Hospital rule, we have
tan x sec 2 x
lim
x 0
lim 1.
x x0 1
[CH – 2017 (1 mark)]: Let i and j be the unit vectors in the x and y directions, respectively. For
the function F ( x, y ) x 3 y 2 , the gradient of the function, i.e., F is given by
(a) 3 x 2 i 2 y j (b) 6x 2 y (c) 3 x 2 i 2 y j (d) 2 y i 3 x 2 j
3 2 3 2 3 2 3 2
Solution (c): F i j k (x y ) ( x y )i (x y )j ( x y )k
x y z x y z
F 3 x 2 i 2 y j 0 3 x 2 i 2 y j .
x7 2 x5 1
[CS – 2017 (2 marks)]: The value of lim
x 3 3x 2 2 x 1
1
[CS – 2017 (1 mark)]: If f ( x) R sin{ x 2} S , f (1 2) 2 and 0 f ( x)dx 2 R , then the
constant R and S are, respectively
(a) 2 and 16 (b) 2 and 0 (c) 4 and 0 (d) 4 and 16
Solution (c):
f ( x ) R sin{ x 2} S f ( x ) R ( 2) cos{ x 2} f (1 2) 2 R ( 2) cos{ 4}
R4 .
1
1 1 cos{ x 2} R
Also 0 f ( x) dx 2 R [(4 ) sin{ x 2} S ]dx 2 R (4 )
0
2
Sx 2
0
8 8
S 2
2 S 0 . So option (c) is correct.
[EC – 2017 (1 mark)]: Consider the following statements about the linear dependence of the real
valued functions y1 1 , y2 x and y3 x 2 , over the field of real numbers.
I. y1 , y2 and y3 are linearly independent on 1 x 0
II. y1 , y2 and y3 are linearly dependent on 0 x 1
III. y1 , y2 and y3 are linearly independent on 0 x 1
IV. y1 , y2 and y3 are linearly dependent on 1 x 0
[EC – 2017 (2 marks)]: Let I (2 zdx 2 ydy 2 xdz ) where x , y , z are real, and let C be the
C
straight line segment from point A : (0, 2,1) to point B : (4,1, 1) . The value of I is _____.
Solution: The straight line joining point A : (0, 2,1) to point B : (4,1, 1) is given by
x0 y2 z 1
t x 4t , y 2 t , z 1 2t , which is any point on the line joining the
4 0 1 2 1 1
points A and B . So dx 4dt , dy dt , dz 2dt . Also for x 0 , t 0 and for x 4 , t 1 .
t 1
Thus I (2 zdx 2 ydy 2 xdz ) {2(1 2t )(4dt ) 2(2 t )( dt ) 2(4t )( 2dt )}
C t 0
t 1
I {30t 4}dt ( 15t 4t )10 11 .
2
t 0
1 1 x y 1 1 x y
[EC – 2017 (2 mark)]: The values of the integrals 0 0 ( x y )3 dy dx and 0 0 ( x y )3 dx dy
are
(a) same and equal to 0.5 (b) same and equal to –0.5
(c) 0.5 and –0.5, respectively (d) –0.5 and 0.5, respectively
1 1 x y 1 x y 1 2x x y
Solution (c): Let I1 dy dx ; consider I1 3
dy dy
0 3 0 ( x y) 0 ( x y )3
0 ( x y)
1 1
1 1 1 1 ( x y ) 2 ( x y ) 1
I1 2 x dy 0 ( x y )2 dy 2 x 2 1
0 ( x y )3
0 0
1 1 1 1 x 1 1
I1 x 2
2
2
2
( x 1) x ( x 1) x ( x 1) ( x 1) ( x 1)
1
1
1 1 1 1 1 1 x y 1
So I1 dx 1 . Similarly I 2 dx dy .
0 ( x 1) 2 0 0 3
( x 1) 0 2 2 ( x y) 2
[EC – 2017 (2 mark)]: If the vector function F aˆ x (3 y k1z ) aˆ y ( k2 x 2 z ) aˆ z (k3 y z ) is
irrotational, then the values of the constants k1 , k2 and k3 , respectively, are
(a) 0.3, –2.5, 0.5 (b) 0.0, 3.0, 2.0 (c) 0.3, 0.33, 0.5 (d) 4.0, 3.0, 2.0
Solution (b): For F to be irrotational, we must have curl F 0
aˆ x aˆ y aˆ z
x y z 0 aˆ x ( k3 2) aˆ y (0 k1 ) aˆ z ( k2 3) 0
(3 y k1 z ) (k 2 x 2 z ) ( k3 y z )
k1 0, k 2 3, k3 2
[EC – 2017 (2 mark)]: The minimum value of the function f ( x ) (1 3) x ( x 2 3) in the interval
100 x 100 occurs at x _____.
Solution: f ( x ) (1 3) x ( x 2 3) f ( x) (1 3){( x 2 3) x (2 x )} x 2 1 f ( x ) 2 x . For
minimum and maximum, f ( x ) 0 x 1 , so x 1 are the critical points. At x 1, f ( x ) 2
and x 1, f ( x) 2 . So at x 1 , f ( x ) has local minima. Now minimum value of function in the
interval 100 x 100 is min{ f ( 100), f (1), f (100)} min{333433.33, 2 3,333433.33}
333433.33 , which occurs at x 100 .
ex , x 1
[EE – 2017 (2 marks)]: A function f ( x ) is defined as f ( x ) 2
, where x .
ln x ax bx x 1
Which one of the following statement is TRUE?
(a) f ( x ) is NOT differentiable at x 1 for any values of a and b
(b) f ( x ) is differentiable at x 1 for unique values of a and b
(c) f ( x ) is differentiable at x 1 for all values of a and b such that a b e
(d) f ( x ) is differentiable at x 1 for all values of a and b
ex , x 1 ex , x 1
Solution: f ( x) 2
f ( x ) . Now f ( x ) will be
ln x ax bx x 1 (1 x) 2ax b x 1
differentiable at x 1 , if (i) f ( x ) is continuous at x 1 ; (ii) the limit of f ( x ) must exist at x 1 .
(i) f ( x) is continuous at x 1, if LHL x 1 RHL x 1 lim f (1 h) lim f (1 h)
h 0 h 0
1 h 2
lim e lim{ln(1 h) a (1 h) b(1 h)} e a b …(i)
h 0 h 0
2
[EE – 2017 (1 mark)]: Let I c xy dxdy , where R is
R
x 1 2x
y 0
xy 2 dy dx c
5
x 1
x 2x
y 0
y 2 dy dx c
5
x 1
x
3 y0
dx
5
5 8x3 5 8x
4
8 x5 8 55 8 15
I c x dx c x 1 3 dx c c
x 1
3 15 x 1 15 15
I (6 104 ) (24992 15) 0.9996 1 .
x, x 1
1 x, x 0
[EE – 2017 (2 marks)]: Let g ( x ) f ( x) 2 and
. Consider the
x 1, x 1 x , x0
composition of f and g , i.e. ( f g )( x ) f {g ( x )} . The number of discontinuities in ( f g )( x )
present in the interval ( , 0) is:
(a) 0 (b) 1 (c) 2 (d) 4
x, x0 1 x, x0
Solution (a): g ( x) x, 0 x 1 and f ( x) x 2 , 0 x 1. So
x 1, x 1 x2 , x 1
1 ( x ), x0 1 x, x0
2 2
f {g ( x )} ( x ) , 0 x 1 x , 0 x 1 f {g ( x )} is discontinuous at x 0 and
( x 1) ,
2
x 1 ( x 1) ,
2
x 1
x 1 . But in the interval ( , 0) , x 0,1 are not lying; so f {g ( x )} has not discontinuous point in
( , 0) .
2
[MA – 2017 (2 marks)]: Let D be the region in bounded by the parabola y 2 2 x and the line
y x . Then D 3xy dx dy equals _____
x 2 x y
Solution: I 3xy dx dy
D
y 2
y 0 x y
x y2 2 y 0
y2
3 xydx dy 3 y
2 x y2 2
dy
y2
3y
y 2y 4 y2 3 y
3
3 y5 3 y4 3y6 48 3 64
I y2 dy dy 2.
y0 y 0
2 4 2 8 8 48 y 0 8 48
[ME – 2017 (2 marks)]: A parametric curve defined by x cos( u 2) , y sin( u 2) in the range
0 u 1 is rotated about the x axis by 360 degrees. Area of the surface generated is
(a) 2 (b) (c) 2 (d) 4
u dx u u dy u
Solution (c): x cos sin ; y sin cos .
2 du 2 2 2 du 2 2
We know that if the equation of the curve is given in the parametric form x f1 (t ) and y f 2 (t ) , and
the curve revolves about x axis, then the area of surface of revolution
t2 1 dx 2 dy 2
2 f 2 (t ) {( dx dt ) 2 ( dy dt ) 2 }dt . So required surface area A 2 y du
t 0 du
1
du
2 2
u
1 u u 2 1 u
A 2 sin sin cos du 2 sin du
0 0
2 2 2 2 2 2
1
2 2 u
A cos 2
2 0
[ME – 2017 (2 marks)]: The surface integral S F n dS over the surface S of the sphere
div F F i j k {( x y )i ( x z ) j ( y z )k}
x y z
div F ( x y) ( x z) ( y z) 1 0 1 2 .
x y z
By using divergence theorem, S F n dS V (div F)dV , where V is the volume of given surface of
4
sphere x 2 y 2 z 2 9 . So S F n dS V (div F )dV V (2)dV 2 3 (3)
3
226.194 cubic unit.
[MN – 2017 (1 mark)]: The position vector of a moving particle is given by r (t ) t 3 i t j t 2 k .
The acceleration of the particle in the direction of the motion is
(a) 0 (b) 60 i 2 k (c) 6t i 4 j 2 k (d) 6t i 2 k
2
d (r ) d d r d d 3 2 d 2
Solution (d): a 2
(t i t j t k ) (3t i 1 j 2t k ) (6t i 2 k ) .
dt dt dt dt dt dt
(2 x 1) (2 x 1)
[MN – 2017 (1 mark)]: The value of lim ( x 2 2 x 1) (2 x 2 3 x 2)
x
is
[MN – 2017 (2 marks)]: A rectangle has two of its corners on the x axis and the other two on the
parabola y 12 x 2 . The largest area of the rectangle is _____.
Solution: The required area of the rectangle as shown in
figure is given by Area (2 x)(12 x 2 ) 2 x 3 24 x …(i)
For maximum or minimum,
d ( Area)
6 x 2 24 0 x 2, 2 . Also
dx
2 2
d ( Area ) d ( Area )
2
12 x 2
24 0 , so Area is
dx dx x 2
(2 x ) 4 16
[PE – 2017 (1 mark)]: The value of lim is _____.
x 0 x
[PE – 2017 (2 marks)]: For a velocity field given by v y i x j 0 k , calculate the curl of v . If the
calculated vector is a i b j c k , then the value of c is _____.
i j k
Solution: curl(v ) v x y z i (0 0) j (0 0) k ( 1 1) 2 k . So c 2 .
y x 0
2
[PH – 2017 (2 marks)]: The integral 0 x 2 e x dx is equal to _____. (up to two decimal places).
d
Solution:
0
2
0
2
I x 2 e x dx x 2 e x dx x xe x dx
2
( x)
dx
x 2
xe dx dx . 0
Let
2
x2 2 x2 x2 1 t e t e x
I xe dx , put x t 2 xdx dt ; so I xe dx e
2
( xdx ) e dt .
2 2
2 2 2
e e 1
x
ex x
1
Thus I x
2
2
dx x
0 2
0 2
0
e
x2
dx (0 0)
2 0
e
x2
dx
2
The integral 0 e x dx is an Gaussian integral, whose value is (1 2) . So
[PI – 2017 (1 mark)]: Divergence of the curl of a twice differentiable continuous vector function is
(a) unity (b) infinity (c) zero (d) a unit vector
Solution (c): div (curl F) ( F) 0 .
2 t
[PI – 2017 (2 marks)]: The improper integral 0 e dt converges to
(a) 0 (b) 1.0 (c) 0.5 (d) 2.0
2 t 2 20
e e e 1
Solution (c): I e 2t dt 0.5 .
0
2 0 2 2 2
sin y 1
[XE – 2017 (1 mark)]: If 0 x dy dx for some 1 , then the value of is _____.
y 2
sin y sin y
Solution: I x dy dx dx dy , where
0 y D y
D {( x, y) : 0 x ( ) and x y ( )} ; but we can also rewrite D as
D {( x, y ) : 0 x y ( )} or D {( x, y) : 0 y ( ) and 0 x y} . So we have:
sin y y sin y sin y sin y
I
0 x y
dy dx
0 0 y
dx dy
0 y
0
y
dx dy
0 y y dy
1
I sin y dy ( cos y )0 cos 1 . As I ; so
0 2
1 1
cos 1 cos 3 .
2 2 3
[XE – 2017 (2 marks)]: Let f : and g : be defined by
x (sin x ) cos(1 x ), x 0 x cos(1 x), x 0
f ( x) and g ( x) , where denotes the set of real
0, x0 0, x0
numbers. Then, at x 0 ,
(a) f is differentiable but g is NOT differentiable
(b) f is not differentiable but g is differentiable
(c) both f and g are differentiable
(d) neither f nor g is differentiable
f (0 h) f (0) (0 h){sin(0 h)}cos{1 (0 h)} 0
Solution: f (0 ) lim lim
h 0 h h 0 h
f (0 ) lim{ sin h}cos{1 h} 0 (a number between 1 and 1) 0
h 0
g (0 ) a number between 1 and 1 , so g (0 ) does not exist. Hence no need to find g (0 ) , as
k 2
[XE – 2017 (2 marks)]: Let ak 2 k k 4 sin k and bk 2 k sin 2 k for k 1, 2, . Then
(a) k 1 ak converges but k 1 bk does NOT converges
(b) k 1 ak does NOT converge but k 1 bk converges
(c) both k 1 ak and k 1 bk converge
(d) neither k 1 ak nor k 1 bk converges
Solution (b): Since sin k 1 2 k k 4 sin k 2 k k 4 ak 2 k k 4 ; so k 1 ak k 1 2 k k 4 .
Now using the ratio test for k 1 2 k k 4 , we have, where ak 2 k k 4 ,
ak 1 2 ( k 1) ( k 1) 4 {1 (1 k ) 4 1
L lim
k a
lim
k k 4
2 k
lim
k 2
1.
2
So k 1 2 k k 4 converges; thus
k
k 1 ak also converges as k 1 ak k 1 2 k k 4 . Also we know that if k 1 ak converges then
k 1 ak converges; hence k 1 ak , i.e. k 1 2 k k 4 sin k also converges.
2 k2 2 k 2 k 2 2
k 2
Now sin k 1 2 k sin k 2 k k 1 2 k sin k k 1 2 k .
2 2
So using ratio test for k 1 2 k k , where bk 2k k we have
2
bk 1 2 ( k 1) ( k 1) {1 (1 k ) {1 (1 k ) 2
L lim
k b
lim
k k2
lim 2
k ( k 1) k
2 lim
k 22 k 1
0 1. Thus k 1 2 k k
k 2 k 2
2 2 2
converges; hence k 1 2 k k sin 2 k also converges as k 1 2 k k sin 2 k k 1 2 k k .
So option (b) is correct.
Chapter 2 : Calculus
2.1 Set Theory
The symbols used in defining the sets are given as:
Symbol Meaning Symbol Meaning
Implies There exist
Belongs to iff If and only if
A B A is a subset of B & And
Implies and is implied by a|b a is a divisor of b
Does not belong to N Set of Natural numbers
s.t. Such that I or Z Set of Integers
For every R Set of real numbers
A set is a well-defined collection (i.e. a rule which decides whether a given element belongs or
does not belong to the given collection) of elements. A set is described in the following two ways:
In Roster method, a set is described by listing elements, in any order, separated by commas,
within braces {} . The set of vowels of English alphabet is represented as {a, e, i, o, u} .
In Set-builder method, a set is described by a characterizing property P ( x ) of its elements x as
{x : P( x ) holds} . The set E of all even natural numbers can be written as E { x : x 2n, x N } .
Types of Sets
Null set: A set which contains no element at all is called ‘null or empty or void set’. It is denoted
by the symbol or {} . A set which has at least one element is called a non-empty set.
Singleton set: A set consisting of a single element is called a singleton set.
Finite set: A set is called a finite set if it is either void set or its elements can be counted by
natural number, and the process of counting terminates at a certain natural number n (say), which
is called as the cardinal number or order of a finite set A and is denoted by n( A) or O ( A) .
Infinite set: A set whose number of elements is infinite or cannot be counted is called infinite set.
Equivalent set: Two finite sets A and B are equivalent if their cardinal numbers are same i.e.
n( A) n( B ) . For e.g., A {1,3,5} ; B {6,9,11} are equivalent sets [ O( A) O( B ) 4] .
Equal set: Two sets A and B are said to be equal, i.e. A B if x A x B . Equal sets are
always equivalent but equivalent sets may need not be equal set.
Universal set: A universal set is a set, denoted by U , containing of all possible elements which
occur in the discussion. The universal set is not unique; it may differ from problem to problem.
Power set: If S is any set, then the family of all the subsets of S , i.e. P( S ) {T : T S } , is
called the power set of S . The power set of S is denoted by P( S ) . Obviously and S are both
elements of P( S ) . Let S {a, b, c} , then P ( S ) {},{a},{b},{c},{a, b},{a, c},{b, c},{a, b, c} .
If set A , then P( A) has one element n[ P ( A)] 1 , i.e. P ( A) is always non-empty.
n
If set A has n elements, then P( A) has 2 elements.
P P ( ) ,{ } P P P ( ) ,{ }, { } , ,{ } . Hence n P P P ( ) 4 .
Subsets: Let A and B be two sets. If every element of A is an element of B , then A is called a
subset of B , and we write A B . Thus, A B a A a B . The total number of subset of a
n
finite set containing n elements is 2 .
If A is a subset of B and A B then A is a proper subset of B . We write this as A B .
The null set is subset of every set and every set is subset of itself, i.e., A and A A for
every set A , they are called improper subsets of A .
Disjoint sets: Two sets A and B are said to be disjoint, if sets A and B have no element in
common, i.e. if A B , then A and B are called disjoint sets, as shown in Fig. 2.3.
Difference of sets: Let A and B be two sets. The difference A B , is the set of all those elements
of A which do not belong to B . Thus, A B {x : x A and x B} , as shown in Fig. 2.4, the
shaded part represents A B .
Symmetric difference of two sets: Let A and B be two sets. The symmetric difference of sets
A and B is the set ( A B ) ( B A) and is denoted by A B , as shown in Fig. 2.5. Thus,
A B ( A B) ( B A) { x : x A B} .
Complement of a set: Let U be the universal set and let A be a set such that A U . Then, the
c
complement of A with respect to U is denoted by A or A or U A and is defined the set of all
those elements of U which are not in A , as shown in Fig. 2.6. Thus, Ac {x U : x A} .
Figure 2.1: Union of Sets A and B Figure 2.2: Intersection of Sets and Figure 2.3: Disjoint Sets and
Some Important Results on Number of Elements in Sets: If A , B and C are finite sets and
U be the finite universal set, then
n( A B ) n( A) n( B ) n( A B )
If A and B are disjoint non-void sets then n( A B ) n( A) n( B ) .
n( A B) n( A) n( A B ) n( A B ) n( A B ) n( A)
n( A B) n ( A B ) ( B A) n( A B ) n( B A) [ ( A B) & ( B A) are disjoint]
n( A B ) n( A) n( A B) n( B ) n( B A) n( A) n( B) 2n( A B )
n( A B C ) n( A) n( B ) n(C ) n( A B ) n( B C ) n(C A) n( A B C )
n( Ac B c ) n( A B )c n(U ) n( A B ) n( Ac B c ) n( A B )c n(U ) n( A B )
Example 2.2 [CS-1998 (2 marks)]: In a room containing 28 people, there are 18 people who speak
English, 15 people who speak Hindi and 22 people who speak Kannada. 9 persons speak both English
and Hindi, 11 persons speak both Hindi and Kannada whereas 13 persons speak both Kannada and
English. How many people speak all three languages?
(a) 9 (b) 8 (c) 7 (d) 6
Solution (d): n( E ) 18 , n( H ) 15 , n( K ) 22 , n( E H ) 9 , n( H K ) 11 , n( K E ) 13 ,
and n( E H K ) 28 . We have to find n( E H K ) . We know that
n( E H K ) n( E ) n( H ) n( K ) n( E H ) n( H K ) n( K E ) n( E H K ) . So
28 18 15 22 9 11 13 n( E H K ) n( E H K ) 6 .
Example 2.4 [ME-2006 (1 mark)]: Let x denote a real number. Find out the INCORRECT
statement.
(a) S x : x 3 represents the set of all real numbers greater than 3
(b) S x : x 2 0 represents the empty set
(c) S x : x A and x B represents the union of set A and set B
(d) S x : a x b represents the set of all real numbers between a and b , where a and b are
real numbers
Solution (c): A B {x : x A or x B} .
c c
Example 2.5 [CS-1996 (1 mark)]: Let A and B be sets and let A and B denote the complements
of the sets A and B . The set ( A B ) ( B A) ( A B ) is equal to
(a) A B c c (c) A B c c
(b) A B (d) A B
Solution (a): As shown in figure the total shaded region shows
( A B ) ( B A) ( A B ) which is equal to A B .
Example 2.6 [CS-2005 (1 mark)]: Let A , B and C be non-empty sets and let X ( A B ) C and
Y ( A C ) ( B C ) . Which of the following is true?
(a) X Y (b) X Y (c) Y X (d) None of these
Solution (a): X ( A B ) C X ( A C ) ( B C ) Y
Example 2.7 [CS-2006 (2 marks)]: Let E , F and G be finite sets. Let X ( E F ) ( F G ) and
Y E ( E G ) E F . Which one of the following is true?
(a) X Y (b) X Y (c) X Y (d) X Y and Y X
Solution (c): From the Venn diagrams shown below, we can say that X Y .
2.1.2 Relations
Let a relation R from set A to set B is a subset of the Cartesian product of A and B , i.e., A B .
Thus R is a relation from A to B R A B . If R is a relation from a non-empty set A to a
non-empty set B and if ( a, b) R , then we write aRb which is read as ‘ a is related to b by the
relation R ’. If ( a, b) R , then we write a R b which is read as ‘ a is not related to b by the relation
R ’. For e.g., if A {1, 2,3} and B {a, b, c} , then R {(1, b), (2, c ), (3, a )} being a subset of A B ,
is a relation from A to B . Here (1, b) , (2, c) and (3, a) R , so we write 1Rb , 2Rc , 3Ra . But
(2, b) R so we write 2 R b . A relation from set A to set B can be represented in following ways:
(a) In Roster form, a relation is represented by the set of all ordered pairs belonging to R . The
relation R from set A {1, 2} to B {1, 2, 3, 4} by the rule aRb a 2 b is R {(1,1), (2, 4)} .
(b) In Set builder form, the relation R from set A to B is represented as
R {( a, b) : a A, b B and satisfy the rule which associates a and b} . If A {1, 2,3} and
B {1,1 2 ,1 3,1 4} and R is a relation from A to B given by
R {(1,1), (2,1 2), (3,1 3)} , then in set builder form R can be described as
R {( a, b) : a A, b B and b 1 a} . It should be noted that it is not possible
to express every relation from set A to set B in set builder form. (c) In arrow
diagram, a relation from set A to a set B is represented by drawing arrows
from first components to the second components of all ordered pairs belonging
to R. If A {a, b, c} and B { p, q} , then a relation Figure 2.8: Arrow
diagram representation
R {( a, p ), (a, q ), (b, q ), (c, p )} , shown in Fig. 2.8.
of a relation R
Domain and Range of a relation: Let R be a relation from a set A to a
set B . Then the set of all first components of the ordered pairs belonging to R is called the
domain of R , while the set of all second components of the ordered pairs in R is called the range
of R .
Inverse of a relation: Let A and B be two sets and let R be a relation from a set A to a set B .
1
Then the inverse of R , denoted by R , is a relation from B to A and is defined by
1
R 1 {(b, a ) : (a, b) R} . Clearly, ( a, b) R (b, a) R 1 ; also domain of R range of R
1
and range of R domain of R .
Total number of relations: Let A and B be two non-empty finite sets s.t. n( A) m , n( B) n .
mn
Then, A B consists of mn ordered pairs. So total number of subsets of A B is 2 [This
mn
point was asked in CS-1993 (2 marks), CS-1999 (1 mark)]. Among these 2 relations, the
void ( ) and universal ( A B ) are the trivial relations from A to B .
Types of Relations
Void Relation: A relation R on a set A is called void or empty relation, if no element of A is
related to any element of A . For e.g., a relation R on the set A {1, 2,3} is defined by
R {(a, b) : a b 10} ; as a b 10 for any two elements of A , so ( a, b) R for any a, b A
and thus R does not contain any element of A A R is an empty set.
Universal relation: A relation R on a set is called universal relation if each element of A is
related to every element of A . For e.g., a relation R on the set A {1, 2,3} is defined by
R {( a, b) R : ( a b) 0} ; as a b 0 for all a, b A , so ( a, b) R for all ( a, b) A A ,
i.e., each element of set A is related to every element of set A and thus R A A R is a
universal relation on set A .
Identity relation: Let A be a set. Then the relation I A {( a, a) : a A} on A is called the
identity relation on A , i.e., every element of A is related to itself only. For e.g., if A {1, 2,3}
the R1 {(1,1), (2, 2), (3,3)} is an identity relation; but R2 {(1,1), (2, 2)} and
R3 {(1,1), (2, 2), (3,3), (1, 2)} are not identity relation as (3,3) R2 and in R3 1 is related to 2.
Example 2.8: Three relations R1 , R2 and R3 are defined on set A {a, b, c} as follows: (i)
R1 {( a, a ), ( a, b), ( a, c ), (b, b), (b, c ), (c, a), (c, b), (c, c)} (ii) R2 {( a, b), (b, a ), ( a, c ), (c, a )} (iii)
R3 {(a, b), (b, c), (c, a )} . Find whether R1 , R2 and R3 is reflexive, symmetric and transitive.
Solution: For (i): R1 {( a, a ), ( a, b), ( a, c ), (b, b), (b, c ), (c, a), (c, b), (c, c)} . Clearly,
( a, a ), (b, b), (c, c ) R1 R1 is reflexive on A . As, ( a, b) R1 but (b, a ) R1 R1 is not symmetric
on A . As, (b, c) R1 and (c, a ) R1 but (b, a ) R1 R1 is not transitive on A .
For (ii): R2 {( a, b), (b, a ), ( a, c ), (c, a )} . Clearly, ( a, a), (b, b), (c, c ) R2 R2 is not reflexive on A
. As, ( a, b), (b, a ) R2 and ( a, c ), (c, a ) R2 so R2 is symmetric on A . As, ( a, b) R2 and
(b, a) R2 but ( a, a ) R2 R2 is not transitive on A .
For (iii): R3 {(a, b), (b, c), (c, a )} . Clearly, ( a, a ), (b, b), (c, c) R3 R3 is not reflexive on A . As,
(b, c ) R3 but (c, b) R3 R3 is not symmetric on A . As, ( a, b) R3 and (b, c ) R3 but
( a, c ) R3 R3 is not transitive on A .
Example 2.9 [CS-1994 (1 mark)]: Amongst the properties {reflexivity, symmetry, anti-symmetry,
transitivity} the relation R {( x, y ) N 2 x y} satisfies ……….
Solution: As x y ( x, x ) R R is not reflexive relation. If ( x, y ) R then ( y , x) R for all
x, y N R is symmetric relation. If ( x1 , y1 ) R and ( y1 , x1 ) R for all x1 , y1 N then for R to
be anti-symmetric we must have x1 y1 which contradicts x1 y1 and so R is not anti-symmetric. If
( x1 , y1 ) R and ( y1 , x1 ) R for all x1 , y1 N then for R to be transitive ( x1 , x1 ) R which
contradicts x1 x1 so R is not transitive.
Example 2.10 [CS-1995 (1 mark)]: Let R be a symmetric and transitive relation on a set A . Then
(a) R is reflexive and hence equivalence relation (b) R is reflexive and hence partial order
(c) R is reflexive and hence not an equivalence (d) None of these
relation
Solution (d): If R be a symmetric and transitive relation on a set A , then it is not necessary that R is
reflexive. Let us consider an example, let a relation R {( a, b), (b, a ), ( a, a), (b, b)} on the set
A {a, b, c} , as R is symmetric and transitive both but as (c, c ) R R is not reflexive.
Example 2.11 [CS-1996 (2 marks)]: Let R be a non-empty relation on a collection of sets defined
by A R B if and only if A B . Then, (pick the true statement)
(a) R is reflexive and transitive (b) R is symmetric and not transitive
(c) R is an equivalence relation (d) R is not reflexive and not symmetric
Solution (d): As there is no common elements in set A and B so any element of A is not related to
same element and thus R is not reflexive. So option (a) and (c) are not correct. If ( a, b) R then for
R to be symmetric (b, a) R set A and set B have a common element ‘ b ’ which contradicts the
statement A B thus R is not symmetric. So option (b) is wrong and option (d) is correct.
Example 2.12 [CS-1998 (2 marks)]: The binary relation R {(1,1), (2,1), (2, 2), (2, 3), (2, 4), (3,1),
(3, 2), (3,3), (3, 4)} on the set A 1, 2,3, 4 is
(a) reflexive, symmetric and transitive (b) neither reflexive, nor irreflexive but transitive
(c) irreflexive, symmetric and transitive (d) irreflexive and antisymmetric
Solution (b): As (4, 4) R R is not reflexive relation; so option (a) is wrong. As (2,1) R but
(1, 2) R R is not symmetric relation, so option (c) is wrong. As (1,1), (2, 2), (3,3) R R is not
irreflexive relation, so option (d) is wrong. As for every ordered pair if ( a, b) R and (b, c ) R then
( a, c ) R and thus R is transitive relation. So the given relation is neither reflexive, nor irreflexive
but transitive.
Example 2.13 [CS-2002 (2 marks)]: A binary relation S (empty set) on set A {1, 2,3} is
Exercise: 2.1
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
10. Sets A and B have 3 and 6 elements respectively. What can be the minimum number of
elements in A B ? _____
11. Let n(U ) 700, n( A) 200, n( B ) 300 and n( A B ) 100, then n( Ac B c ) _____
12. In a town of 10,000 families it was found that 40% family buy newspaper A , 20% buy newspaper
B and 10% families buy newspaper C , 5% families buy A and B , 3% buy B and C and 4%
buy A and C . If 2% families buy all the three newspapers, then number of families which buy
A only is _____.
13. In a city 20% of the population travels by car, 50% travels by bus and 10% travels by both car and
bus. Then persons (in %) travelling by car or bus is _____.
14. Suppose A1 , A2 , A3 , , A30 are thirty sets each having 5 elements and B1 , B2 , Bn are n sets each
30 n
with 3 elements. Let Ai B j S and each elements of S belongs to exactly 10 of the Ai ’s
i 1 j 1
29. Out of 800 boys in a school, 224 played cricket, 240 played hockey and 336 played basketball. Of
the total, 64 played both basketball and hockey; 80 played cricket and basketball and 40 played
cricket and hockey; 24 played all the three games. The number of boys who did not play any
game is _____.
30. A survey shows that 63% of the Americans like cheese whereas 76% like apples. If x % of the
Americans like both cheese and apples, then
(a) x 39 (b) x 63 (c) 39 x 63 (d) None of these
31. Of the members of three athletic teams in a school 21 are in the cricket team, 26 are in the hockey
team and 29 are in the football team. Among them, 14 play hockey and cricket, 15 play hockey
and football, and 12 play football and cricket. Eight play all the three games. The total number of
members in the three athletic teams is _____.
32. In a college of 300 students, every student reads 5 newspaper and every newspaper is read by 60
students. The number of newspaper is
(a) at least 30 (b) at most 20 (c) exactly 25 (d) None of these
33. In the village called Sultankari, only three TV channels are available: Soon Plus, Mony and See
TV. Out of 4000 TV viewers in the village, 1500 watch Soon TV, 2000 watch Mony and 2500
watch See TV. Amongst these, 500 viewers watch Soon Plus and Mony, 800 watch Soon Plus and
See TV, and 1000 watch Mony and See TV. How many viewers watch all three channels? _____
34. There is a shortage of tube-lights, bulbs and fans in a village Sultankari. All houses do not have
either tube-light or bulb or fan. Exactly 19% of houses do not have just one of these. Atleast 67%
of houses do not have tube-lights. Atleast 83% of houses do not have bulbs. Atleast 73% of
houses do not have fans. What percentage (in%) of houses does not have tube-light, bulb and fan?
_____
35. If A is a non-empty subset of a set E , then what is E ( A ) ( A ) ?
(a) A (b) A
c (c) (d) E
36. In a class of 110 students, x students take both Mathematics and Statistics, 2 x 20 students take
Mathematics and 2 x 30 students take Statistics. There are no students who take neither
Mathematics nor Statistics. What is the value of x ? _____
37. In a school there are 30 teachers who teach Mathematics or Physics. Of these teachers, 20 teach
Mathematics and 15 teach Physics, 5 teach both Mathematics and Physics. The number of
teachers teaching only Mathematics is _____.
38. Which one of the following is correct?
(a) {} {{},{{ }}} (b) {} {{ },{{ }}} (c) {{},{{ }}} (d) {{},{{}}}
39. If A {(22 n 3n 1) | n N } and B {9(n 1) | n N } then which one of the following is
correct?
(a) A B (b) B A (c) A B (d) Neither (a) nor (b) is correct
c
40. {( A B ) A} ( A B ) is (a) (b) A (c) B (d) B
c
(a) Only (i) (b) Only (ii) (c) Only (iii) (d) (ii) and (iii)
46. If A {1, 2, 3, 4} , then what is the number of subsets of A with atleast three elements? _____
47. Which of the following is correct statement?
(a) (b) P ( ) (c) P( ) (d) P ( )
48. Let R be the relation on the set R of all real numbers defined by aRb iff a b 1 . Then R is
(a) Reflexive and Symmetric (b) Symmetric only
(c) Transitive only (d) Anti-symmetric only
49. Consider the following statements: (i) Set of points of a given line is a finite set; (ii) Intelligent
students in a class is a set; (iii) Good books in a school library is a set. Which of the above
statements is/are not correct?
(a) Only (i) (b) Both (ii) and (iii) (c) Both (i) and (ii) (d) (i), (ii) and (iii)
50. In an examination, 52% candidates failed in English and 42% failed in Mathematics. If 17%
candidates failed in both English and Mathematics, what percentage (in %) of candidates passed
in both the subjects? _____
51. Which one of the following is a correct statement?
(a) {a} {{a},{b},{c}} (b) {a} {{a}, b, c} (c) {a, b} {{a}, b, c} (d) a {{a}, b, c}
52. Let X {1, 2,3, 4, 5} and Y {1,3, 5, 7,9} . Which one of the following is not a relations from X
to Y ?
(a) R1 {( x, y ) | y 2 x, x X , y Y } (b) R2 {(1,1), (2,1), (3, 3), (4,3), (5, 5)}
(c) R3 {(1,1), (1,3)(3,5), (3, 7), (5, 7)} (d) R4 {(1,3), (2, 5), (2, 4), (7, 9)}
53. The relation ‘less than’ in the set of natural numbers is
(a) Only symmetric (b) Only transitive (c) Only reflexive (d) Equivalence relation
54. With reference to a universal set, the inclusion of a subset in another, is relation, which is
(a) Symmetric only (b) Equivalence relation (c) Reflexive only (d) None of these
55. Let A {2, 4, 6,8} . A relation R on A is defined by R {(2, 4), (4, 2), (4, 6), (6, 4)} . Then R is
(a) Anti-symmetric (b) Reflexive (c) Symmetric (d) Transitive
56. Let N denote the set of all natural numbers and R be the relation on N N defined by
( a, b) R (c, d ) if ad (b c ) bc ( a d ) , then R is
(a) Symmetric only (b) Reflexive only (c) Transitive only (d) An equivalence relation
57. In order that a relation R defined on a non-empty set (a) is reflextive (b) is symmetric
A is an equivalence relation, it is sufficient, if R (c) is transitive (d) is equivalence
58. Let R and S be two relations on a set A . Then which one of the following is not correct?
(a) R and S are transitive, then R S is also transitive
(b) R and S are transitive, then R S is also transitive
(c) R and S are reflexive, then R S is also reflexive
(d) R and S are symmetric then R S is also symmetric
2.2 Function
A function is a relation whose every input corresponds with a
single output. Every function is a relation; but every
relation is not a function. Let X and Y be any two non-
empty sets. “A function from X to Y is a rule that assigns to
each element of set X , one and only one element of set Y ”.
Let the rule be ‘ f ’ then mathematically we write f : X Y
where y f ( x ), x X and y Y . We say that ‘ y ’ is the
image of ‘ x ’ under f (or x is the pre image of y ). Two
things should always be kept in mind:
A mapping f : X Y is said to be a function if each
element in the set X has its image in set Y (as shown in
Fig. 2.9). It is also possible that there are few elements in Figure 2.9: A mapping f:X→Y
set Y which are not the images of any element in set X .
Every element in set X should have one and only one image. Functions cannot be multi-valued
(A mapping that is multi-valued is called a relation from X and Y ).
Testing for a function by vertical line test: A relation
f : A B is a function or not it can be checked by a graph
of the relation. If a vertical line is drawn which cuts the
given curve at more than one point, as shown in Fig. 2.10
(i) and (ii), then the given relation is not a function; and
when it cuts the curve at only one point, as shown in Fig.
2.10 (iii) and (iv), then the given relation is a function
Number of functions: Let X and Y be two finite sets Figure 2.10: Vertical line test
having n and m elements respectively. Then each element
of set X can be associated to any one of m elements of set Y . So, total number of functions
from set X to set Y is m n . [This point was asked in CS-1998 (1 mark)]
Domain, Co-domain and Range of Function: As shown in Fig. 2.11, if a function f is defined
from set A to set B then for f : A B set A is called the domain of function f and set B is called
the co-domain of function f . The set of all f images of the elements of A is called the range of
function f . In other words, Domain All possible values of x for which f ( x ) exists; Range For
all values of x , all possible values of f ( x ) ; Range Co-domain.
Methods for finding domain of a function: The expressions under even root must be (i.e.,
square root, fourth root etc.) 0 ; the
denominator must be 0 . If domain of
y f ( x) and y g ( x) are D1 and D2
respectively then the domain of f ( x ) g ( x ) or
Figure 2.11: Domain, Co-domain and Range of a
f ( x ) g ( x ) is D1 D2 ; while domain of
Function
f ( x) g ( x) is D1 D2 {g ( x) 0} .
Methods for finding Range of a function: Range of y f ( x) is collection of all outputs f ( x )
corresponding to each real number in the domain. If domain finite number of points range
set of corresponding f ( x ) values. If domain or [some finite points], then express x
in terms of y ; from this we find y for x to be defined (i.e., find the values of y for which x
exists). If domain a finite interval, find the least and greatest value for range using
monotonicity. There is no specific method to find out the range of a function, as the type of
method varies with different types of functions.
Two function f and g are said to be equal functions, if and only if (i) domain of f domain
of g ; (ii) Co-domain of f co-domain of g ; (iii) f ( x) g ( x ) x their common domain.
Example 2.17 [PI-2013, ME-2013 (1 mark)]: Choose the CORRECT set of functions, which are
linearly dependent
(a) sin x , sin 2 x and cos 2 x (b) cos x , sin x and tan x
2 2
(c) cos 2x , sin x and cos x (d) cos 2x , sin x and cos x
Solution (c): cos 2 x cos x sin x , so cos 2x , sin 2 x and cos 2 x are linearly dependent.
2 2
Example 2.20 [EC-2007 (1 mark)]: Which one of the following function is strictly bounded?
2
(a) 1 x 2 (b) e x (c) x 2 (d) e x
Solution (d): For y 1 x 2 , as x 0, y is y 1 x2 unbounded. For y e x , as
x , y y e x is unbounded. For y x 2 , as x , y y x 2 is unbounded. For
2 2
y e x , as x 0, y 1 and x , y 0 , also any value of x ( , ) , y e x lies between
2
y (0,1] ; and thus y e x is bounded.
Please note that the questions on ‘Logarithms’ also came in previous years ‘General Aptitude’
section in GATE examination. These questions and exercise on it are given in ‘Section 9.4 of
General Aptitude’.
Explicit and implicit functions: A function is said to be
explicit if it can be expressed directly in terms of the independent
variable. If the function cannot be expressed directly in terms of
the independent variable or variables, then the function is said to
be implicit. e.g. y sin 1 x log x is explicit function, while
x 2 y 2 xy and x 3 y 2 ( a x ) 2 (b y ) 2 are implicit functions.
Constant function: Let k be a fixed real number. Then a
function f ( x ) given by f ( x) k for all x R is called a
constant function. The domain of the constant function
f ( x) k is the complete set of real numbers and the range of
f is the singleton set {k} . The graph of a constant function
is a straight line parallel to x axis as shown in Fig. 2.13 (a)
and it is above or below the x axis according as k is
positive or negative. If k 0 , then the straight line coincides
with x axis.
Identity function: The function defined by f ( x ) x for
all x R , is called the identity function on . Clearly, the Figure 2.13: (a) Constant (b) Identity (c)
domain and range of the identity function is . The graph of Modulus Function
the identity function is a straight line passing through the origin and inclined at an angle of 45o
with positive direction of x axis, as shown in Fig. 2.13 (b).
Number of onto function (surjection): The number of onto functions from set A (having m
n
elements) to set B (having n elements), where 1 n m , is r 1 (1)n r nCr r m .
Any polynomial function f : R R is, Onto function if its degree is odd; Into function if its
degree is even.
Example 2.26 [CS-2012 (2 marks)]: How many onto (or surjective) functions are there from an n
element ( n 2) set to a 2 element set?
(a) 2
n n
(b) 2 1
n
(c) 2 2 (d) 2(2n 2)
Solution (c): The number of onto functions from set A (containing n elements) to set B (containing
2
2 elements) is r 1 (1)2 r 2Cr r n (1)21 2C11n (1)22 2C2 2n 2 2 n .
One-One and Onto function (Bijective): A function f : A B is said to be a Bijective function
if it is one-one as well as onto.
Number of One – One and Onto function (Bijective): If A and B are finite sets and
f : A B is a bijection, then A and B have the same number of elements. If A has n
elements, then the number of bijection from A to B is the total number of arrangements of n
items taken all at a time i.e., n! .
Example 2.28 [CS-1996 (2 marks)]: Let R denotes the set of real numbers. Let f : R R R R
be a bijective function defined by f ( x, y ) ( x y , x y ) . The inverse function of f is given by
1
(a) f ( x, y) 1 ( x y ) ,1 ( x y) (b) f 1 ( x, y ) ( x y ) 2 , ( x y ) 2
1 1
(c) f ( x, y ) x y, x y (d) f ( x, y) 2( x y ), 2( x y)
Solution (b): As the given function is bijection so its inverse exists f 1 f ( x, y ) f 1 ( x y, x y )
f 1 ( x y , x y ) ( x, y ) . Now let x y X and x y Y , on solving these two equations we
get, x (X Y) 2 and y (X Y) 2 . Thus,
f 1 ( x y, x y ) ( x, y ) f 1 ( X , Y ) ( X Y ) 2, ( X Y ) 2 . Now putting x in place of X and
y in place of Y , we get f 1 ( x, y ) ( x y ) 2 , ( x y ) 2 .
Periodic Function: A function is said to be periodic function if its each value is repeated after a
definite interval. So a function f ( x ) will be periodic if a positive real number T exist such that,
f ( x T ) f ( x ), x domain. Here the least positive value of T is called the period of the function.
Clearly f ( x ) f ( x T ) f ( x 2T ) f ( x 3T ) ..... . sin n x , cos n x , sec n x and cos ec n x has
n n
period, if n is even; 2 if n is odd or fraction. tan x, cot x has period if n is odd or even.
sin x , cos x , tan x , cot x , sec x and cosec x are periodic with period of . The following are
the rules for finding the period of periodic function:
If f ( x ) is periodic with period T , then af ( x ) b , a ( 0), b R , is also periodic with period T .
If f ( x ) has a period T , then the function f ( ax b ) , a ( 0), b R , will have a period T a .
If f ( x ) is periodic with period T then 1 f ( x) , f ( x ) is also periodic with same period T .
Let f ( x ) has period p m n ( m, n N and co-prime) and g ( x ) has period q r s ( r , s N
and co-prime) and let t be the LCM of p and q , then t LCM( m, r ) HCF( n, s) . LCM of
rational and irrational number is not possible. If m and r are irrational then LCM( m, r ) does not
exist unless they have same irrational surd.
Let us suppose that f ( x ) is periodic with period p and g ( x ) is periodic with period q . Let r be
the LCM of p and q , if it exists. Then if f ( x ) and g ( x ) cannot be interchanged by adding a
least positive number k r , to x , then r is the period of f ( x ) g ( x) . If f ( x ) and g ( x ) can be
interchanged by adding a least positive number k r , to x , then k is the period of f ( x ) g ( x) .
A constant function is periodic but its period is not defined. [This point was asked in AE-2011].
Example 2.30 [EC-2009 (1 mark)]: A function given by f (t ) sin 2 t cos 2t . Which of the
following is true? f has frequency components at
(a) 0 and 1 2 Hz (b) 0 and 1 Hz (c) 1 2 and 1 Hz (d) 0, 1 2 and 1 Hz
2
Solution (b): f (t ) sin t cos 2t (1 cos 2t ) 2 cos 2t (1 cos 2t ) 2 . As cos t is of period 2
cos 2t is periodic with period 2 2 ; so (1 cos 2t ) 2 is periodic with period . So the time
period ( T ) of the given function is and so its frequency is 1 T 1
Example 2.31 [IN-2009 (1 mark)]: The fundamental period of x (t ) 2 sin 2 t 3sin 3 t , with t
expressed in second, is
(a) 1 (b) 0.67 (c) 2 (d) 3
Solution (d): period of 2sin 2 t is 2 2 1 ; period of 3sin 3 t is 3 2 3 2 ; so period of
x (t ) is LCM(1 1, 3 2) LCM(1,3) HCF(1, 2) 3 1 3
Example 2.32 [EE-2010 (1 mark)]: The period of the signal x (t ) 8 sin 0.8 t 4 is
(a) 0.4 s (b) 0.8 s (c) 1.25 s (d) 2.5 s
Solution (d): The period of the signal x (t ) 8 sin 0.8 t 4 is same as the period of
x(t ) sin(0.8 t ) . As period of sin t is 2 period of sin(0.8 t ) is 2 0.8 2.5
such that g (b) y . SO we have shown that for every y C , there is a b B such that g (b) y
, so g is onto.
Exercise: 2.2
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
27. The function f ( x ) sin log( x x 2 1 ) is
(a) Even function (b) Odd function (c) Neither even nor odd (d) Periodic function
28. The function f ( x) sin( x 2) 2 cos( x 3) tan( x 4) is periodic with period _____.
29. If f ( x ) is an odd periodic function with period 2, then f (4) _____.
30. The period of f ( x ) x [ x ] is _____.
31. If a, b be two fixed positive integers such that (a) a (b) b
3 2 2 3 13
f ( a x ) b [b 1 3b f ( x ) 3b{ f ( x )} { f ( x)} ] for all real x (c) 2a (d) 2b
then f ( x ) is a periodic function with period
32. f ( x ) sin 2 x sin 2 x ( 3) cos x cos x ( 3) and g (5 4) 1, then ( gof ) ( x) _____
33. If g ( x ) x 2 x 2 and (1 2)( gof )( x ) 2 x 2 5 x 2, then f 1 ( x )
(a) 2 x 3 (b) 2 x 3 (c) 2 x 2 3 x 1 (d) 2 x 2 3 x 1
34. Let g ( x ) 1 x [ x ] , where [] denotes greatest integer function, and f ( x ) sgn( x ) , then for
all x , f ( g ( x )) _____
35. If the function f : R R be such that f ( x ) x [ x ] where [ y ] denotes the greatest integer less
than or equal to y , then f 1 ( x ) is
(a) 1 ( x [ x]) (b) [ x ] x (c) Not defined (d) None of these
x ( x 1) 1
36. If f : [1, ) [1, ) is defined as f ( x ) 2 then f ( x ) is equal to
x ( x1)
(a) (1 2) (b) (1 2)(1 1 4 log 2 x ) (c) (1 2)(1 1 4 log 2 x ) (d) Not defined
37. Which one of the following function is invertible?
(a) f ( x) 2 x (b) f ( x ) x 3 x (c) f ( x) x 2 (d) None of these
Left Hand Limit (LHL) and Right Hand Limit (RHL): Consider the values of the functions
at the points which are very near to a on the left of a . If these values tend to a definite unique
number as x tends to a , then the unique number so obtained is called left-hand limit (LHL) of f ( x )
at x a and symbolically we write it as f ( a 0) lim f ( x ) lim f ( a h) , here we write a h in
x a h 0
place of x a ; where h is very-very small positive number. Similarly we can define right-hand limit
(RHL) of f ( x ) at x a which is expressed as f ( a 0) lim f ( x ) lim f ( a h) , here we write
x a h 0
a h in place of x a ; here h is very-very small positive number. If both LHL and RHL exists and
both are equal, i.e., lim f (a h) lim f (a h) , then we say that lim f ( x ) exists.
h 0 h 0 x a
Example 2.33 [CH-2009 (2 marks)]: The value of the limit lim cos x {x ( 2)}3 is
x 2
sin( h)
RHL 1 lim 1 {(any values between 0 and 1) } 1 0 1 . So, RHL RHL 1
( h )
h 0
Algebra of limits: The following hold true for lim f ( x ) l and lim g ( x ) m ( l , m R ).
x 0 x 0
lim( f ( x ) g ( x )) l m lim( f ( x ) g ( x )) l m
x a x a
lim{ f ( x ) g ( x )} l m lim k f ( x) kl
x a x a
lim g ( x )
lim log{ f ( x)} log{lim f ( x )} lim[ f ( x )]g ( x ) {lim f ( x )}xa
x a x a
x a x a
Some important expansions: In finding limits, use of expansions of some functions are useful:
(1 x ) n 1 nx
n( n 1) 2 ( x log a) 2
x a x 1 x log a
2! 2!
x x2 x3 x3 x5
e 1 x [This series was sin x x [This series was
2! 3! 3! 5!
asked in EC-1995, CE-2012 (1 mark)] asked in CE-1998 (1 mark)]
x 2 x 4 x6 x3 2 x5
cos x 1 tan x x
2! 4! 6! 3 15
[This series was asked in CE-1997 (1 x2 x3 x 4
mark)] log(1 x ) x , x 1
2 3 4
x2 x3 x4 x 3
x 5
log(1 x ) x , x 1 sin h x x
2 3 4 3! 5!
x2 x4 x6 x3 5
cosh x 1 tanh x x 2 x
2! 4! 6! 3
1 2 x3 2 2 x5 1 x3 x5 x7
sin x x 1 3 1 tan x x
3! 5! 3 5 7
2
x x
1x (1 x ) log(1 x )
1
x 11
(1 x) e e 2 3
e 1 x2
2 24
Example 2.36 [ME-2007 (2 marks)]: lim e 1 x ( x 2 2)
x 0
x
x3
(a) 0 (b) 1/6 (c) 1/3 (d) 1
Example 2.37 [CE-2014 (2 marks)]: The expression lim ( x 1) is equal to
0
(a) log x (b) 0 (c) x log x (d)
Solution (a): In the given limit, is the variable; using expansion of x by treating x as constant
and being variable, we put in place of x , x in place of a in expansion of a x ,
x 1 ( log x )
2
(log x ) 2
lim lim 1 log x 1 lim log x log x
0 0
2! 0
2!
Evaluation of Algebraic limits: The following are the methods for evaluating algebraic limits:
Direct substitution method: If by direct substitution of the point in the given expression we get a
finite number, then the number obtained is the limit of the given expression.
Example 2.38 [MN-2009 (1 mark)]: The value of lim 2 4 x 2 5 is
x2
[Similar questions were also asked in ME-2000, CE-2004, TF-2008, (1 marks), TF-2010 (2
marks)]
1/3
Example 2.40 [ME-2008 (1 mark)]: The value of lim ( x 2) ( x 8) is
x 8
(a) 1/16 (b) 1/12 (c) 1/8 (d) ¼
1/3 1/3
x 2 ( x 2) 1 1
Solution (b): lim 1/3 3 3
lim 1/3 2/3 1/3 2
lim 2/3 1/3 2
x 8 ( x ) 2 x 8 ( x 2)( x 2 x 2 ) x8 x 2 x 2 12
[Similar question was also asked in MN-2010 (1 mark)]
Based on the form when x : In this case, expression should be expressed as a function of
1 x and then after removing indeterminate form, (if it is there) replace 1 x by 0. Step I: Write
down the expression in the form of rational function, i.e. f ( x) g ( x ) , if it is not so. Step II: If k
is the highest power of x in numerator and denominator both, then divide each term of numerator
n
and denominator by x k . Step III : Use the result lim (1 x ) 0 , where n 0 .
x
( x 1) x 1 (1 x ) 1 0 1
L lim lim
2 2
x
( x x 1 x) x x
1 (1 x) (1 x ) 1) 1 0 0 1) 2
n( n 1) n( n 1)(n 2)
Proof: Since, (1 x) n 1 nx x2 x 3 , where x 1 , n Q (from
2! 3!
binomial expansion). When x is infinitely small (approaching to zero) such that we can ignore
n
higher power of x , then we have (1 x) 1 nx .
LHL lim
n
x a
n
lim
n
( a h) a
n
lim
a n (1 ( h a )) n 1 lim a 1 (nh a ) 1 na
n
n 1
x a xa h 0 aha h 0 h h 0 h
LHL lim
n
x a
n
lim
n
( a h) a
n
lim
an (1 ( h a )) n 1 lim a 1 (nh a) 1 na
n
n 1
x a xa
aha h0 h h 0 h0 h
n n
x a n 1
Hence, if n Q , then lim na
x a x a
n
Example 2.47 [TF-2011 (1 mark)]: The value of lim{(1 x) 1} x is
x 0
(a) 0 (b) n (c) (d) 1 n
Solution (b): Using binomial expansion, the given limit becomes
(1 x ) n 1 n( n 1) 2 n( n 1)(n 2) 3
lim 1 nx x x 1 x
x 0 x 2! 3!
n( n 1) n(n 1)(n 2)
lim n x x2 n
x 0 2! 3!
Example 2.48: Find the values of constants a and b so that lim
x
( x 2
1) ( x 1) ax b 0
Solution: lim
x
( x 2
1) ( x 1) ax b 0 lim x
x 2
(1 a) x ( a b) 1 b ( x 1) 0
Since the limit of the given expression is zero, therefore degree of the polynomial in numerator must
be less than that of denominator. As the denominator is a first degree polynomial. So, numerator must
be a constant i.e., a zero degree polynomial. 1 a 0 and a b 0 a 1 and b 1 .
Evaluation of Trigonometric limits: The following results holds true for f ( x) 0 when x a .
sin f ( x) sin x lim cos f ( x ) 1 lim cos x 1
lim 1 lim 1 [This point x a x 0
x a f ( x) x 0 x
tan f ( x ) tan x
was asked in IN-2008, ME-2011, MT-2009 lim 1 lim 1
(1 mark), CS-1995 (2.5 marks)]
x a f ( x) x0 x
sin 1 f ( x ) sin 1 x sin x cos x
lim 1 lim 1 lim lim 0
x a f ( x) x 0 x x x x x
lim sin 1 x sin 1 a, a 1 lim cos 1 x cos 1 a; a 1
x a x a
1 1
tan f ( x) tan x
lim 1 lim 1 lim tan 1 x tan 1 a; a
x a
x a f ( x) x 0 x
sinf ( x )
Proof of lim 1 : Using expansion of sin x , we have
x a f ( x)
lim
sin f ( x )
lim f ( x)
f ( x)3 f ( x)5 f ( x)7
f ( x)
xa f ( x)
xa 3! 5! 7!
sin f ( x ) f ( x )2 f ( x)4 f ( x )6
lim lim 1 1 , since as x a , f ( x) 0 .
xa f ( x) x a
3! 5! 7!
The other results can also be proved in a similar manner.
Example 2.49 [ME-1994 (1 mark)]: The value of lim 1 sin x 1 tan x is
x 0
(a) 0 (b) 1 (c) 2 (d)
2
1 1 1 cos x 2 sin ( x 2)
Solution (a): lim lim lim lim tan( x 2) 0
x 0 sin x tan x x 0 sin x x 0
2 sin( x 2) cos( x 2) x0
Example 2.50 [ME-1995 (2 mark)]: Evaluate lim ( x 3 cos x ) ( x 2 sin 2 x)
x
(a) (b) 0 (c) 2 (d) Does not exist
Example 2.56 [BT-2014 (1 mark)]: The limit of the function e 2t sin(t ) as t , is …………
2 t
Solution: lim e sin(t ) lim e 2t lim sin(t ) 0 (any value between 1 to 1) 0
t t t
sin[ x]
, [ x] 0
Example 2.57: If f ( x ) [ x ] , then find lim f ( x ) .
x 0
0 , [ x] 0
Solution: In closed interval of x 0 at right hand side [ x ] 0 and at left hand side [ x ] 1. Also
sin[ x]
, ( 1 x 0)
[0] 0. Therefore function is defined as f ( x) [ x ]
0 , (0 x 1)
sin[ x ] sin( 1)
Left hand limit lim f ( x ) lim sin1c
[ x]
x 0 1 x 0
sin( x 4)
Example 2.58 [CS-2016 (1 mark)]: lim _____
x 4 x4
sin( x 4) sin( x 4) sin h
Solution: lim lim lim 1
x 4 x4 x 4 0 x4 h 0 h
Evaluation of Exponential and Logarithmic limits based on series expansion: The following
standard forms are used in case f ( x) 0 , when x a
1 f ( x) 1x x
lim 1 f ( x) e lim 1 x e or lim 1 1 x e [This limit was asked in EC-
x a x 0 x
2014 (1 mark)]
b f ( x) 1 ax 1 ex 1
lim log e b ( b 0 ) lim log e a or lim 1
xa f ( x) x 0 x x 0 x
log 1 f ( x) log 1 x
lim 1 lim 1
x a f ( x) x0 x
1 f ( x) 1 f (x)
Proof of lim 1 f ( x ) e : Using binomial expansion we have, lim 1 f ( x)
x a x a
1 1 1 2 1 1 1 3
f ( x) 1 f ( x ) 1 2 f ( x )
f ( x) f ( x) f ( x ) f ( x ) f ( x ) f ( x)
lim 1
x a 1! 2! 3!
1 f (x) 1 1 f ( x ) 1 f ( x ) 1 2 f ( x ) 1 1
lim 1 f ( x) lim 1 1 1 e
xa x a
1! 2! 3! 2! 3!
The other forms can be proved in a similar manner.
c dx d
lim 1 1 (a bx)
x
a bx
e and lim
x a bx
b
log e (1 4 x )
Example 2.62 [ME-2016 (1 mark)]: lim is equal to
x 0 e3 x 1
(a) 0 (b) 1 12 (c) 4 3 (d) 1
x
log(1 x ) e 1
Solution (c): We know that lim 1 ; thus 1 and lim
x x 0 x 0 x
4 x[{log e (1 4 x )} 4 x ] 4 lim[{log e (1 4 x)} 4 x ] 4
x 0
L lim 3 x
3 x
.
x 0 3x{(e 1) 3 x} 3 lim[{(e 1) 3 x}] 3
x0
[Similar question was also asked in PI-2016 (2 marks)]
L’Hospital’s rule: If f ( x) and g ( x) be two functions of x such that (i) lim f ( x ) lim g ( x ) 0 ;
x a x a
(ii) Both are continuous at x a ; (iii) Both are differentiable at x a ; (iv) f ( x ) and g ( x ) are
continuous at the point x a , then lim f ( x ) g ( x ) lim f ( x ) g ( x ) provided that g ( a ) 0 .
x a xa
putting x 0 the given expression take the form 0 0 ; so again differentiating Nr and Dr separately
w.r.t. x , we get lim (sin x ) (cos x ) 0 1 0
x 0
ex 1
Example 2.67 [PE-2016 (1 mark)]: The value of lim is equal to _____.
x 0
sin x
Solution: At x 0 the given expression takes 0 0 form, so using L’Hospital rule, i.e. differentiating
ex 1 ex e0 1
Nr. and Dr. separately w.r.t. ‘ x ’, we get, L lim lim 1.
x 0 sin x
x0 cos x cos 0 1
[Similar question was also asked in ME-2014 (1 mark)]
Sandwich theorem: Sandwich theorem helps in calculating the limits, when limits cannot be
calculated using the above discussed methods. According to this theorem, if f ( x ) , g ( x ) and h( x) are
any three functions such that, f ( x ) g ( x ) h( x ) x neighbourhood of x a and
lim f ( x) lim h( x) l ( say ) , then lim g ( x ) l. This theorem is normally applied when the lim g ( x)
x a xa x a x a
can't be obtained by using conventional methods as function f ( x ) and h( x) can be easily found.
Example 2.68: Evaluate lim (log x) [ x]
x
Example 2.69 [CE-2011 (2 marks)]: What should be the value of such that the function
cos x ( 2) x if x 2
f ( x) defined below is continuous at x 2 ?
1 if x 2
(a) 0 (b) 2 (c) 1 (d) 2
Solution (c): The turning points for f ( x ) is x 2 and f ( 2) 1 (given). So,
cos( 2 h) sinh
LHL at x 2 lim f ( 2 h) lim lim
h 0 h 0 ( 2) ( 2 h ) h 0 h
cos( 2 h) sinh
RHL at x 2 lim f ( 2 h) lim lim
h 0 h 0 ( 2) ( 2 h) h 0 h
If f ( x ) is continuous at x 2 , then LHL RHL f ( 2) 1
Example 2.70 [CS-2013 (1 mark)]: Which one of the following functions is continuous at x 3 ?
2, x3
4, x3
(a) f ( x ) x 1, x3 (b) f ( x )
( x 3) 3, x 3 8 x , x 3
x 3, x3 1
(c) f ( x ) (d) f ( x ) 3
, x3
x 4, x 3 x 27
Solution (a): For option (a), LHL at ‘ x 3 ’ lim f (3 h) lim (3 h) 3 3 2 ; RHL at ‘ x 3 ’
h 0 h 0
1 1
lim lim 3 , so at x 3 , LHL RHL ; hence f ( x ) is not
h 0 (3 h) 3 27 h 0 h 9 h 3h 2
continuous at x 3 .
Example 2.71 [MN-2014 (2 marks)]: The value of a , for which the function below is continuous at
2 x ax 2 , x 1
x 1 is f ( x )
4 x 3,x 1
(a) –5 (b) 0 (c) 5 (d) 10
Solution (c): The turning points for f ( x ) is x 1 and f (1) 2 1 a 1 2 a . So,
2
h 0 h 0
LHL at x 1 lim f (1 h) lim 2(1 h) a(1 h) 2 lim(2 2 h a 2ha ah 2 ) 2 a
h 0
Example 2.72 [ME-2016 (1 mark)]: The values of x for which Solution (c): The given function is
x 2 3x 4 not continuous for all x where
the function f ( x ) 2 is NOT continuous are denominator term is zero, i.e.
x 3x 4
(a) 4 and –1 (b) 4 and 1 x 2 3x 4 0
(c) –4 and 1 (d) –4 and –1 ( x 4)( x 1) 0 x 4,1 .
[ x ] [ x ], x (0, 2) {1}
we redefine f ( x ) as, f ( x ) . Now f ( x ) is continuous for x (0, 2) .
1, x 1
Limit and Continuity of a Function of two variables: When we extend the continuity of a
function of one variable to functions of two variables, we will see that the domain of functions of two
2 2
variables is a subset of , in other words it is a set of pairs. A point in is of the form ( x , y ) . So,
the equivalent of x a will be ( x, y ) ( a, b) . For functions of three variables, the equivalent of
x a will be ( x, y, z ) (a, b, c) , and so on. x could only approach a from two directions, from
the left or from the right, but ( x , y ) can approach ( a, b) from infinitely many directions. In fact, it
does not even have to approach ( a, b) along a straight path as shown in Fig. 2.18. For functions of
several variables, we would have to show that the limit along every possible path exists and are
the same. The problem is that there are infinitely many such paths. To show a
limit does not exist, it is still enough to find two paths along which the limits
are not equal. In view of the number of possible paths, it is not always easy to
know which paths to try. We can try some suggested following paths: (i)
Horizontal line through ( a, b) , the equation of such a path is y b ; (ii)
Vertical line through ( a, b) , the equation of such a path is x a ; (iii) Any
straight line through ( a, b) , the equation of the line with slope m through Figure 2.18: Limit of a
( a, b) is y mx b am ; (iv) Quadratic paths; for e.g., a typical quadratic function in two
variables
paths through (0, 0) is y x 2 .
Properties of Limits of Functions of two (or more) variables: Let us assume that L , M and k are
real numbers and that lim f ( x, y ) L and lim g ( x, y ) M , then
( x , y ) ( a ,b ) ( x , y ) ( a , b )
Factorization method: In this method, numerator and denominator are factorised. The common
factors are cancelled and the rest outputs the results. For e.g.,
3 3 2 2
x y ( x y )( x xy y )
lim lim lim ( x 2 xy y 2 ) 0
( x , y ) (0,0) x y ( x , y ) (0,0) x y ( x , y ) (0,0)
Limit along a path: We usually use this method to prove the limit of a function in two variable
does not exists.
( x 4 xy )
Example 2.75 [XE-2008 (1 mark)]: is lim
( x3 y 3 ) ( x , y ) (0,0)
along one path exists, but along another path it does not exists. So the given limit does not exists.
xy
Example 2.76 [CE-2016 (1 mark)]: What is the value of lim 2 2
?
x 0, y 0 x y
(a) 1 (b) 1 (c) 0 (d) Limit does not exist
Solution (d): Limit along the path x 0 : First we find what the function becomes along this path
xy (0) y xy
2 2
2 2
0 lim 2 2
lim (0) 0 .
x y x 0 (0) y x 0, y 0 x y x 0, y 0
Limit along the path y x : First we find what the function becomes along this path
xy x2 1 xy 1 1
2 2
2
lim 2 2
lim .
x y yx
2x 2 x 0, y 0 x y x 0, y 0 2 2
As limit along the path x 0 is to the limit along the path y x and so given limit does not exist.
Change of coordinate: Sometimes changing coordinates may be useful in evaluating the limits of
two variables. For e.g., putting x r cos and y r sin in lim
( x , y ) (0,0)
(x 3
y 3 ) ( x 2 y 2 ) , the
limit ( x, y ) (0, 0) changes to r 0 and thus we get
3
lim r (cos sin )
r 0
3 3
r 2 2 2
(cos sin ) lim r (cos sin ) 0
r 0
3 3
Sandwich Theorem: This theorem is also useful in evaluating the limits of two variables. For
e.g., in evaluating the lim ( x 2 y ) ( x 2 y 2 ) , we need to find a function g ( x, y ) such that
( x , y ) (0,0)
2 2
x2 y x2 y x y x y
0 f ( x, y ) 0 g ( x , y ) ; thus f ( x, y ) 0 0 y
x2 y2 x2 y2 x2 y 2 x2 y 2
2 2 2 x2 y x2 y
(as x ( x y ) ) lim 0 lim 2 2
lim y 0 lim 2 2
0
( x , y ) (0,0) ( x , y ) (0,0) x y ( x , y ) (0,0) ( x , y ) (0,0) x y
lim ( x2 y) ( x 2 y 2 ) 0 lim ( x2 y) ( x 2 y 2 ) 0
( x , y ) (0,0) ( x , y ) (0,0)
Similarly, slope of the tangent at point P , which is the limiting position of the chords drawn on the
right hand side of point P , which is denoted as ‘Right Hand Derivative (RHD),’ i.e.,
RHD lim f ( a h) f (a ) h lim slope of chord PR lim slope of chord PR
h 0 h 0 R P
If h 0 , then, f (c h ) f (c ) f (c h ) f (c ) f (c ) f (c h ) f (c ) h h
lim f (c h ) lim f (c ) lim f (c h ) f (c) h lim h lim f ( c) f (c) 0 f (c )
h 0 h 0 h 0 h 0 h 0
Similar argument with one – sided limits show that if f has a derivative from one side (right or
left) at x c , then f is continuous from that side at x c .
If f ( x ) is differentiable at x a LHD and RHD both exists and LHD RHD There is a
unique tangent (which happens when the curve is smooth) at point P .
A function whose graph is not smooth will fail to have a derivative where the graph has
a corner, where the one sided derivative differ, as shown in Fig. 2.20(a)
a cusp, where the slope of PQ
approaches from one side and
from the other, as shown in
Fig. 2.20(b).
a vertical tangent, where the slope
of PQ approaches from both
sides or approaches from both
Figure 2.20: Types of Discontinuities
sides, as shown in Fig. 2.20(c)
a discontinuity, as shown in Fig. 2.20(d).
A function f ( x ) defined in an open interval ( a, b) is said to be differentiable or derivable in open
interval ( a, b) if it is differentiable at each point of ( a, b) .
A function f : [ a, b] R is said to be differentiable in [ a, b] if
f ( x ) exists for every x such that a x b i.e., f is differentiable in ( a, b)
Right hand derivative of f at x a exists
Left hand derivative of f at x b exists
If f ( x ) and g ( x ) both are not differentiable at x a , then the sum function may be a
differentiable function.
Example 2.77 [ME-2002 (1 mark)]: Which of the following function is not differentiable in the
domain [ 1,1] ?
(a) f ( x) x 2 (b) f ( x ) x 1 (c) f ( x ) 2 (d) f ( x ) max( x, x )
Solution (d): As the graph for the function in option (d) has a sharp turn at x 0 , so it is not
differentiable at x 0 . All other functions in options (a), (b) and (c) have no sharp turn for all x ,
so they are differentiable for all x .
Example 2.78 [AE-2008 (1 mark)]: The (a) is neither continuous nor differentiable at
sin x x 0 x0
(b) is continuous nor differentiable at x 0
function defined by f ( x ) 0 x0
(c) is differentiable but not continuous at x 0
3x2 x 0 (d) is continuous but not differentiable at x 0
Solution (d): For continuity and differentiability of f ( x ) , the turning point of f ( x ) is x 0 . As
f (0) 0 and LHL x 0 lim f (0 h) lim sin(0 h) lim( sinh) 0 ;
h 0 h 0 h 0
2 2
RHL x 0 lim f (0 h) lim 3(0 h) lim 3h 0
h 0 h0 h 0
RHD x 0 lim f (0 h) f (0) h lim 3(0 h) 0
h 0 h0
2
h lim (3h 2 ) h lim 3h 0
h 0 h 0
LHD x 0 lim
h 0
h (h 3 5
3!) ( h 5!) h 2
h 0
3 5
( h ) lim ( h 3!) ( h 5!) ( h 2 )
h h3
LHD x 0 lim 0 . Similarly,
h 0
3! 5!
f (0 h) f (0) sin(0 h) (0 h) 1 sinh h
RHD x 0 lim lim lim
h 0 h h0 h h 0 h2
LHD x 0 lim
h ( h3 3!) ( h5 5!) h
lim
( h3 3!) ( h5 5!)
h 0 h2 h 0 h2
h h3
LHD x 0 lim 0 . LHD x 0 RHD x 0 0 f ( x ) is differentiable at x 0 .
h0
3! 5!
sin(1 x ), x 0
Example 2.81 [AE-2014 (1 mark)]: The function given by f ( x ) is
0, x0
(a) unbounded everywhere (b) bounded and continuous everywhere
(c) bounded but not continuous at x 0 (d) continuous and differentiable everywhere
Solution (c): As ‘ sin ’ function is always bounded whose values lies between [ 1,1] . For the given
function the turning point is x 0 . So
LHL x 0 lim f (0 h) lim sin 1 (0 h) limsin(1 h) ; as 1 h sin(1 h) lies between
h 0 h 0 h 0
[ 1,1] ; so lim sin(1 h) has oscillating limit and thus it does not exists. Similarly,
h 0
[ 1,1] ; so lim sin(1 h) has oscillating limit and thus it does not exists. Hence the given function is not
h 0
continuous at x 0 .
D (cosec 1 x ) 1 x
x 2 1 , for x 1 D (sin h 1 x ) 1 1 x2
D (sin h x ) cos h x D (cos h 1 x) 1 x2 1
D (cos h x ) sin h x D (tan h 1 x ) 1 ( x 2 1)
D (tan h x ) sec h 2 x D (cot h 1 x ) 1 (1 x 2 )
D (cot h x) cosec h 2 x
D (sec h 1 x) 1 x 1 x 2
D (sec h x ) sec h x tan h x
D (cosec h 1 x ) 1 x 1 x 2
3
Example 2.84 [IN-2007 (2 marks)]: Consider the function f ( x ) x , where x is real. Then the
function f ( x ) at x 0 is
(a) continuous but not differentiable (b) once differentiable but not twice
(c) twice differentiable but not thrice (d) thrice differentiable
x, x 0 x 3 , x 0
3
Solution (c): x 0, x 0 f ( x ) x 0, x0.
x, x 0 x3 , x 0
3 3
As, LHL x 0 lim f (0 h) lim{(0 h) } 0 ; RHL x 0 lim f (0 h) lim{(0 h) } 0 ; and
h 0 h0 h0 h 0
Example 2.86 [EC-2016 (1 mark)]: Given the following statements about a function f : ,
select the right option:
P: If f ( x ) is continuous at x x0 , then it is also differentiable at x x0 .
Q: If f ( x ) is continuous at x x0 , then it may not be differentiable at x x0 .
R: If f ( x ) is differentiable at x x0 , then it is also continuous at x x0 .
(a) P is true, Q is false, R is false (b) P is false, Q is true, R is true
(c) P is false, Q is true, R is false (d) P is true, Q is false, R is true
Solution: We know that if a function f ( x ) is differentiable at x x0 , then it is also continuous at
x x0 . But if a function f ( x ) is continuous at x x0 , then f ( x ) may or may not be differentiable at
x x0 . Thus statements Q and R are true, and P is false; so option (b) is correct.
1 11 h
1
f (0 h) f (0) 1 (0 h) 1
LHD x 0 lim lim lim 1 h lim 1 . Similarly
h 0 h h 0 h h 0 h h 0 1 h
1 1 1 h
1
f (0 h) f (0) 1 (0 h) 1
RHD x 0 lim lim lim 1 h lim 1
h0 h h0 h h 0 h h0 1 h
As LHD x 0 RHD x 0 , so the given function is not differentiable at x 0 . So option (b) is correct.
x2 a2 x a tan or a cot (a 2 x 2 ) (a 2 x2 )
2 2
x a cos 2
x 2 a2 x a sec or a cos ec ( x a)( x b ) 2
x a sec b tan
2
2
ax x 2 x a sin ( x a)(b x ) 2
x a cos b sin
2
2 2
x ( a x) x a tan x ( a x) x a sin
Quotient rule:
d f ( x)
dx g ( x ) d
dx
d
dx 2
g ( x ) ( f ( x )) f ( x ) ( g ( x )) ( g ( x )) , provided g ( x ) 0
Example 2.89 [ME-2005 (2 marks)]: With a 1 unit change in b , what is the change in x in the
solution of the system of equation x y 2 , 1.01x 0.99 y b ?
(a) Zero (b) 2 units (c) 50 units (d) 100 units
Solution (c): Solving the given two equations, we get 0.02 x b 0.02dx db
dx db 0.02 50(db) 50 units (as db 1 unit is given).
Differentiation of infinite series: If y is given in the form of infinite series of x and we have to
find out dy dx then we remove one or more terms, it does not affect the series.
2
If y f ( x) f ( x) f ( x ) , then y f ( x ) y y f ( x) y
dy dy dy f ( x )
2y f ( x )
dx dx dx 2 y 1
f ( x ) f ( x ) y
f (x)
If y f ( x) then y f ( x) log y y log f ( x)
1 dy y f ( x ) dy dy y 2 f ( x )
log f ( x )
y dx f ( x) dx dxf ( x )[1 y log f ( x )]
1 1 dy yf ( x )
If y f ( x ) y f ( x)
1 y dx 2 y f ( x )
f ( x)
f ( x) ....
Example 2.98: Find the differential coefficient of tan 1 2 x (1 x 2 ) w.r.t. sin 1 2 x (1 x 2 ) .
Solution: Let y1 tan 1 2 x (1 x 2 ) and y2 sin 1 2 x (1 x 2 ) . Putting x tan
1 1 2
y1 tan tan 2 2 2 tan x dy1 dx 2 (1 x ) , and
y2 sin 1 sin 2 2 tan 1 x dy2 dx 2 (1 x 2 ) . Hence dy1 dy2 (dy1 dx) (dy2 dx ) 1
Example 2.101 [CS-2014 (2 marks)]: The function f ( x) x sin x satisfies the following equation:
f ( x ) f ( x) t cos x 0 . The value of t is _____
Solution: f ( x ) x sin x f ( x ) sin x x cos x f ( x ) 2 cos x x sin x . So
f ( x ) f ( x) t cos x 0 2 cos x x sin x x sin x t cos x 0 (t 2) cos x 0 t 2 .
Example 2.102: If the volume of a spherical balloon is increasing at the rate of 900 cm2/sec, then find
the rate of change of radius of balloon at instant when radius is 15 cm [in cm/sec]
Solution: V (4 3) r 3 , Differentiate with respect to t , we get
dV 4 dr dr 1 dV dr 1 1 7
3r 2 2
900 .
dt 3 dt dt 4 r dt dt 4 15 15 22
Exercise: 2.3
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. lim
x a
a 2 x 3x
3a x 2 x equals
5. lim
x 2
1 tan( x 2) [1 sin x] 1 tan( x 2) [ 2 x]3
(a) 1 8 (b) 0 (c) 1 32 (d)
[( a n) nx tan x]sin nx
6. If lim 0, where n is non-zero real number, then a is equal to
x 0 x2
(a) 0 (b) ( n 1) n (c) n (d) n (1 n)
7. lim log e (1 2h) 2 log e (1 h) h 2
h 0
17. Given that f (2) 6 and f (1) 4 , then lim f (2h 2 h 2 ) f (2)
h 0
f (h h 2
1) f (1)
_____.
x a 2 2 sin x , 0 x 4
18. If the function f ( x ) x cot x b , 4 x 2 , is continuous in the interval [0, ]
b sin 2 x a cos 2 x , 2 x
then the values of ( a, b) are
(a) ( 1, 1) (b) (0, 0) (c) ( 1,1) (d) (1, 1)
19. The function f ( x ) log(1 ax ) log(1 bx ) x is not defined at x 0 . The value which should
be assigned to f at x 0 so that it is continuous at x 0 , is
(a) a b (b) a b (c) log a log b (d) log a log b
20. Let f ( x ) be defined for all x 0 and be continuous. Let f ( x ) satisfy f ( x y) f ( x) f ( y ) for
all x, y and f (e) 1, then
(a) f ( x ) ln x (b) f ( x ) is bounded
(c) f (1 x) 0 as x 0 (d) x f ( x) 1 as x 0
21. The function f ( x ) x x x is
(a) Continuous at the origin
(b) Discontinuous at the origin because x is discontinuous there
(c) Discontinuous at the origin because x x is discontinuous there
(d) Discontinuous at the origin because both x and x x are discontinuous there
22. If the function f is defined by f ( x ) x (1 x ) , then at what points f is differentiable
(a) Everywhere (b) Except at x 1 (c) Except at x 0 (d) Except at x 0 or 1
23. The value of the derivative of x 1 x 3 at x 2 is _____.
(sin x 2 ) x , x 0
24. The function f defined by f ( x )
0 , x0
(a) Continuous and derivable at x 0 (b) Neither continuous nor derivable at x 0
(c) Continuous but not derivable at x 0 (d) None of these
25. If f ( x ) log x , then
(a) f ( x ) is continuous and differentiable for all x in its domain
(b) f ( x ) is continuous for all x in its domain but not differentiable at x 1
(c) f ( x ) is neither continuous nor differentiable at x 1
(d) None of these
26. The function f ( x) x x 1 is
(a) Continuous at x 1 , but not differentiable (b) Both continuous and differentiable at x 1
(c) Not continuous at x 1 (d) None of these
ax 2 b , x 1
27. If the derivative of the function f ( x ) 2 is everywhere continuous and
bx ax 4 , x 1
differentiable at x 1 then
(a) a 2, b 3 (b) a 3, b 2 (c) a 2, b 3 (d) a 3, b 2
28. Let f be twice differentiable function such that f ( x ) f ( x ) and f ( x) g ( x ),
h( x ) { f ( x )}2 {g ( x )}2 . If h(5) 11, then h(10) is equal to _____.
2 x 3 [ x] , x 1
29. The function f ( x )
sin( x 2) , x 1
(a) Is continuous at x 2 (b) Is differentiable at x 1
(c) Is continuous but not differentiable at x 1 (d) None of these
x2
30. The set of points where the function f ( x) 1 e is differentiable
(a) ( , ) (b) ( , 0) (0, ) (c) ( 1, ) (d) None of these
x
31. The function f ( x ) e is
(a) Continuous everywhere but not differentiable at x 0 (b) Not continuous at x 0
(c) Continuous and differentiable everywhere (d) None of these
2
32. If f ( x) 1 1 x , then f ( x ) is
(a) Continuous on [ 1,1] and differentiable on ( 1,1)
(b) Continuous on [ 1,1] and differentiable on ( 1, 0) (0,1)
(c) Continuous and differentiable on [ 1,1]
(d) None of these
33. The function f ( x ) ( x 2 1) x 2 3 x 2 cos x is not differentiable at x _____.
xe (1 x ) (1 x )
, x0
34. If f ( x ) , then f ( x ) is
0 , x0
(a) Continuous as well as differentiable for all x
(b) Continuous for all x but not differentiable at x 0
(c) Neither differentiable nor continuous at x 0
(d) Discontinuous every where
35. Let f : R R be a function. Define g : R R by g ( x ) f ( x ) for all x . Then g is
(a) Onto if f is onto (b) One-one if f is one-one
(c) Continuous if f is continuous (d) Differentiable if f is differentiable
dx 2
x
dx
Angle of intersection of two curves: The angle of intersection of two curves is the angle between
the tangents to the two curves at their point of intersection. Let slope of tangent at P ( x1 , y1 ) of
two curves y1 f1 ( x) and y2 f 2 ( x ) is m1 ( dy dx )1( x1 , y1 ) , m2 (dy dx) 2( x1 , y1 ) . As the angle
1
between two straight lines having slopes m1 and m2 is tan (m1 m2 ) (1 m1m2 ) ; so angle
between the tangents of the two curves is found by putting values of m1 and m2 .
Orthogonal curves: If the angle of intersection of two curves is right angle, the two curves
are said to intersect orthogonally. The curves are called orthogonal curve if m1m2 1 .
Length of Tangent, Normal, Sub- Length of tangent
tangent and Sub-normal: As shown
in Fig. 2.21. PA and PB are called
2
PA y cosec y 1 dy dx dy dx
length of tangent and normal, 2
respectively, at point P . If PC be the Length of normal PB y sec y 1 dy dx
perpendicular from P on x axis, Length of sub-tangent
then AC and BC are length of AC y cot y dy dx
subtangent and subnormal,
respectively, at P . Length of sub-normal BC y tan y dy dx
Length of perpendicular from origin to the Tangent: Length of perpendicular from origin
(0, 0) to the tangent at P( x1 , y1 ) of y f ( x) is p y1 x1 (dy dx) ( x1 , y1 ) 1 ( dy dx) 2 .
Example 2.103 [ME-2006 (2 marks)]: Equation of a line normal to f ( x ) ( x 8) 2 3 1 at P (0,5) is
(a) y 3 x 5 (b) y 3 x 5 (c) 3 y x 15 (d) 3 y x 15
Solution (b):
f ( x ) 2 3( x 8)1 3 , so slope of normal at P (0,5)
1 f ( x) 3(0 8)
(0,5)
13
2 3 . Thus, the required equation is y 5 3( x 0) y 3 x 5 .
Example 2.104 [TF-2009 (2 marks)]: Radius of the circle passing through P (2, 4) and the point of
intersection with the x axis of the tangent and normal drawn at P to the curve y 2 8 x is
(a) 1 (b) 2 (c) 3 (d) 4
Solution (d): Slope of tangent at P (2, 4) is
( dy dx ) (2,4) (4 y ) (2,4) 4 4 1 ; Slope of normal at P (2, 4) is
1 ( dy dx ) (2,4) 1 . So eqn of tangent having slope 1 and passing
through P (2, 4) is y 4 1( x 2) y x 2 ; it meets the x axis at
A(2, 0) . Eqn of normal having slope of –1 and passing through P (2, 4)
is y 4 1( x 2) y x 6 ; it meets the x axis at B (6, 0) . Let
the general equation of a circle with centre ( g , f ) and radius
g 2 f 2 c is x 2 y 2 2 gx 2 fy c 0 . As circle passes through P , A and B . So at
P (2, 4) 4 g 8 f c 20 ; at A(2, 0) 4 g c 4 ; at B (6, 0) 12 g c 36 ; on solving,
we get g 2, f 0, c 12 and so radius of the circle is ( 2) 2 (0) 2 ( 12) 4
Example 2.105 [AE-2010 (2 Solution: Any point on the given curve is ( x, e x ) ; slope of a line
marks)]: If e is the base of
the natural logarithms then joining origin to the point ( x, e x ) is (e x 0) ( x 0) e x x …(i);
the equation of the tangent Slope of tangent at ( x, e x ) on the given curve is y ( x ,e x ) e x ex
( x ,e x )
from the origin to the curve
y e x is …(ii). We have e x x e x e x (1 x) 0 x 1 ( e x 0 ) slope
(a) y x (b) y x of a line joining origin to the point ( x, e x ) is e1 1 e . So, the required
(c) y x e (d) y ex equation of tangent is y 0 e( x 0) y ex .
Example 2.106 [CH-2011 (1 mark)]: Solution (d): Slope of lines for the equations given in
Which ONE of the following y ( x ) has the options (a) and (b) are 1 a and a , respectively.
slope of its tangent equal to ax y ? (Note: Differentiating eqn in option (c) w.r.t. x , slope
a and b are real constants)
(a) y ( x b) a (b) y ax b
dy dx x a x2 b x
a y a x ( ay ) .
n
Differentiating eq in option (d) w.r.t. x , slope
(c) y ( x 2 b) a (d) y ax 2 b
dy dx ax ax 2 b ax y .
Example 2.107 [AG- (a) x0 y0 f ( x0 ) , 0 Solution (a): Slope of line joining the points
2012 (1 mark)]: The ( x0 , y0 ) and ( x, 0) is m y0 ( x0 x ) f ( x0 )
tangent line to (b) x0 y0 f ( x0 ) , 0
x x0 y0 f ( x0 ) . So, tangent to
y f ( x) at ( x0 , y0 ) , (c) x0 f ( x0 ) y0 , 0 y f ( x) at ( x0 , y0 ) intersects the x axis at
assuming f ( x) 0 ,
intersects the x axis at
(d) x0 f ( x0 ) y0 , 0 x0 y0 f ( x0 ) , 0 , where f ( x) 0 .
Example 2.108 [CE-2016 (2 marks)]: The angle of intersection (a) 0o (b) 30o
of the curves x 2 4 y and y 2 4 x at point (0, 0) is (c) 45o (d) 90o
Solution (d): Slope of curve C1 : x 2 4 y is given by 2 x 4 y y1 x 2 ( y1 ) (0,0) 0 ; so the
o
tangent to the curve C1 at (0, 0) is making an angle of 0 with the x axis.
Slope of curve C1 : y 2 4 x is given by 2 yy 4 y2 2 y ( y2 ) (0,0) 2 0 undefined ; so the
o
tangent to the curve C 2 at (0, 0) is making an angle of 90 with the x axis. Hence angle between
o o o
the two curves is the angle between the tangent at (0, 0) , which is 90 0 90 .
Example 2.109 [MA-2016 (2 marks)]: Let be the curve which passes through (0,1) and intersects
each curve of the family y cx 2 orthogonally. Then also passes through the point
(a) ( 2, 0) (b) (0, 2) (c) (1,1) (d) ( 1,1)
Solution (a): Let : y f ( x) be a curve which passes through (0,1) , so f (0) 1 …(i). As and
y cx 2 intersect orthogonally, so product of their slopes at intersection point must be 1 .
Since y f ( x) y f ( x) ; and y cx 2 y 2cx y 2 y x ;
So at the intersecting point, we have ( y)(2 y x) 1 ( dy dx)(2 y x) 1 2 ydy xdx …(ii)
On integrating (ii), y 2 ( x 2 2) k , which passes through (0,1) , i.e. (1) 2 (02 2) k k 1 .
Hence the required curve will be y 2 ( x 2 2) 1 , which also passes through ( 2, 0) .
Figure 2.22: Monotonicity of a Function at a point: (a) and (c) are monotonically increasing; (b) and (d) are monotonically
decreasing at =
Necessary and Sufficient condition for Monotonic Function: Let f be a differentiable real
function defined on an open interval ( a, b) .
From Fig. 2.24a, if f ( x ) is strictly increasing function on ( a, b) , then tangent at every point on
the curve y f ( x) makes an acute angle with the ve direction of x axis tan 0
dy dx 0 or f ( x ) 0 x ( a, b) . If there are some points in ( a, b) such that f ( x) 0
then, f ( x) 0 x ( a, b) f ( x ) is monotonically increasing function, as shown in Fig. 2.24b.
From Fig. 2.24c, if f ( x ) is a decreasing function on ( a, b) , then tangent at every point on the
curve y f ( x) makes an obtuse angle with the ve direction of x axis tan 0
dy dx 0 or f ( x ) 0 x ( a, b) . If there are some points in ( a, b) such that f ( x) 0 then,
f ( x) 0 x ( a, b) f ( x ) is monotonically decreasing function, as shown in Fig. 2.24d.
x (t )
Example 2.111 [EE-2006 (2 marks)]: A continuous-time system is described by y(t ) e ,
where, y (t ) is the output and x (t ) is the input. y (t ) is bounded
(a) only when x (t ) is bounded (b) only when x (t ) is non-negative
(c) only for t 0 if x (t ) is bounded for t 0 (d) even when x (t ) is not bounded
x (t )
Solution (d): y (t ) e as x (t ) or , x(t ) x(t ) y (t ) 0 .
Example 2.112: Find the interval in which the function sin 4 x cos 4 x increase in 0, 2 .
4 4 2 2 2 2 2
Solution: f ( x) sin x cos x (sin x cos x) 2sin x cos x (3 4) (1 4) cos 4 x . For f ( x )
to be increasing, f ( x) 0 f ( x ) sin 4 x 0 sin 4 x 0 4 x 3 2 4 x 3 8
at points where f ( x0 ) exists and if f ( x) changes its sign when passing through x x0 then
x0 is called the point of inflection [This point was asked in MT-2011 (1 mark)]. For e.g., for
f ( x) x 3 1 , at x 0 , f ( x) 0 and f ( x ) 0 , so x 0 is the point of inflection [This
example was asked in ME-2012, PI-2012 (1 mark)]. The later condition may be replaced by
f (c ) 0 , when f (c ) exists. Thus, x c is a point of inflection if f (c) 0 and f (c ) 0 .
If f ( x ) 0 , x ( a, b) , then the curve If f ( x ) 0 , x ( a, b) , then the curve
y f ( x) is concave downward in ( a, b) , y f ( x) is concave upward in ( a, b) , as
as shown in Fig. 2.25a. shown in Fig. 2.25b.
Example 2.113 [CE-1999 (1 mark)]: Number of inflection points for the curve y ( x 2) x 4 is
(a) 3 (b) 1 (c) n (d) (n 1) 2
Solution (b): As y ( x 2) x 4 y 5 x 4 8 x 3 y 4 x 2 (5 x 6) y 60 x 2 48 x
For point of inflection, y 0 4 x 2 (5 x 6) 0 x 0, 6 5 . Now y x 0 0 x 0 is not a
point of inflection. Also, y x 6 5 144 5 0 x 6 5 is the point of inflection
Example 2.114 [EC-2016 (1 mark)]: As x varies from –1 to 3 , which one of the following
describes the behaviour of the function f ( x ) x 3 3 x 2 1 ?
(a) f ( x ) increases monotonically (b) f ( x ) increases, then decreases and increases again
(c) f ( x ) increases and then decreases (d) f ( x ) decreases, then increases and decreases again
Solution: f ( x ) x 3 3x 2 1 f ( x ) 3x 2 6 x ; so f ( x) 0 3x( x 2) 0 x 0, 2
So from wavy-curve method, f ( x ) increases in x ( , 0) (2, ) and decreases
in x (0, 2) . So in x ( 1, 0) , f ( x ) increases; then in x (0, 2) , f ( x ) decreases;
and then in x (2,3) , f ( x ) increases. Hence option (b) is correct.
Figure 2.27: First Order Derivative Test for finding Figure 2.28: Second Order Derivative Test for finding
Maxima and Minima point Maxima and Minima point
Second – Order Derivative Test in Ascertaining the Maxima or Minima: As shown in Fig.
2.28, it is clear that as x increases from a h to a h , the function dy dx continuously
decreases, i.e., the function f ( x ) or dy dx continuously decreases which means in the
neighbourhood of x a we have d 2 y dx 2 0 . Similarly, as x increases from b h to b h ,
the function dy dx continuously increases, i.e., the function f ( x ) or dy dx continuously
2 2
increases which means in the neighbourhood of x b we have d y dx 0 . Hence we have:
2 2
At local maxima, i.e., at x a , dy dx 0 and d y dx 0
2 2
At local minima, i.e., at x b , dy dx 0 and d y dx 0
2 2
However, if d y dx 0 , then the test fails. In this case, f can still have a maxima or minima
or point of inflection (neither maxima nor minima), we revert back to the first-order derivative
check for ascertaining the maxima or minima.
Derivative test: It is nothing but the general version of the second – order derivative test. It
n n 1
says that if f (a) f (a) f (a) f (a) 0 and f (a) 0 , then f ( x ) would have a
local maximum or minimum at x a if n is an odd natural number and that x a would be a
n 1 n 1
point of local minima if f (a) 0 and would be a point of local minima if f (a) 0 .
However, if n is even, then f has neither a maximum nor a minimum at x a .
Absolute or Global
Maxima/Minima of a
function in a given
interval: If a function f ( x)
is defined in an interval then
Global maxima/minima is
basically the greatest/least
value of the function in that
interval. Let c1 , c2 ,, cn be
different critical points of the
function y f ( x) in ( a, b) .
Then
Global maxima/minima Figure 2.30: Global Maxima and Minima if a function in open interval
in [ a, b] will be given as,
Global maxima: M1 max{ f (a), f (c1 ), f (c2 ),, f (cn ), f (b)}
Global minima: M 2 min{ f (c1 ), f (c2 ),, f (cn )} , if and only if lim f ( x) M 2 and
x a
Example 2.115 [CE-1999 (5 marks)]: Find the maximum and minimum values of the function
f ( x ) sin x cos 2 x over the range 0 x 2 .
Solution: For the critical points, f ( x ) 0 cos x 2 sin 2 x 0 cos x (1 2 sin x ) 0 cos x 0
or sin x 1 2 . In 0 x 2 , cos x 0 x 2 , 3 2 ; and sin x 1 2 x 6 ,5 6 . So
critical points are x 6 , 2 ,5 6 , 3 2 . As f ( x ) sin x 4 cos 2 x . Thus
f ( x) x 6
5 2 0 x 6 is the point of maxima and thus f ( 6) 1
f ( x) x 2
3 0 x 2 is the point of minima and thus f ( 2) 0
f ( x ) x 5 6
5 2 0 x 5 6 is the point of maxima and thus f (5 6) 1
f ( x) x 3 2
5 0 x 3 2 is the point of minima and thus f (3 2) 2
Thus, local minima of f ( x ) in 0 x 2 occur at x 2 ,3 2 , whose values are 0, –2,
respectively; and local maxima of f ( x ) in 0 x 2 occur at x 6 , 5 6 , whose values are 1, 1,
respectively. Now lim f ( x ) lim f (0 h) lim{sin(0 h) cos 2(0 h)} 1 and
x 0 h0 h 0
lim f ( x ) lim f (2 h) lim{sin(2 h) cos 2(2 h)} 1 . So, we have no global maxima as
x 2 h0 h0
max(1,1) 1 and lim f ( x ) 1 , lim f ( x) 1 ; but we have global minima having value –2, as
x 0 x 2
Example 2.116 [CE-2000 (2 marks)]: The maxima and minima of the function
f ( x ) 2 x3 15 x 2 36 x 10 occur, respectively, at
(a) x 3 and x 2 (b) x 1 and x 3 (c) x 2 and x 3 (d) x 3 and x 4
2
Solution (c): For the critical point of f ( x ) , f ( x) 0 6 x 30 x 36 0 x 2, 3 are the
critical points. As, f ( x ) 12 x 30 , f ( x) x 2 6 0 x 2 is the point of maxima; and
f ( x) x 3 6 0 x 3 is the point of minima.
[Similar questions were also asked in CS-1998 (5 marks), CE-2004, ME-1995, TF-2008, EE-2011
(2 marks), CE-2002 (1 mark)]
(c) (d)
Solution (c): The given symmetric function gives (i) dy dx 0 for x 0 ; and (ii) dy dx 0 for
x 0 . It is to be also noted that, (iii) for x 0 , dy dx increases from zero to ve maximum value
and then decreases to zero; (iv) for x 0 , dy dx decreases from zero to ve minimum value and
then increases to zero. The four conditions are met by the figure given in option (c).
Example 2.119 [EE-2005 (2 marks)]: For the function f ( x ) x 2 e x , the maximum occurs at x
(a) 2 (b) 1 (c) 0 (d) –1
x
Solution (a): For critical point of f ( x ) , f ( x) 0 xe (2 x ) 0 x 0, 2 are the critical
points. Also f ( x ) (2 x )(e x xe x ) xe x f ( x) x 0 2 0 x 0 is the point of minima;
2
and f ( x ) x 2 2e 0 x 2 is the point of maxima.
[Similar question was also asked in EE-2014 (1 marks)]
Example 2.120 [ME-2005 (2 marks)]: The right circular cone of largest volume that can be enclosed
by a sphere of 1 metre radius has a height (in metres) of
(a) 1/3 (b) 2/3 (c) 2 2 3 (d) 4/3
Solution (d): Let ABCA be a sphere of radius OC OA R and ABDCA be a
cone of height AD h and radius CD r . AD AO OD R OD and
from OCD , we have OD OC2 DC2 R2 r 2
AD h R R 2 r 2 1 1 r 2 (as R 1 is given). Volume of the cone
is V (1 3) r 2 h (1 3) r 2 (1 1 r 2 ) . For maximum or minimum volume,
dV r2
0 r 2 2 1 r 2 2
0 2 2 1 r (r
2
1 r 2 ) 0 (as r 0 is not
dr 3 1 r2
possible) or 2 1 r 2 2(1 r 2 ) r 2 2 4 2
2 1 r 2 3r 2 2 4(1 r ) 9r 12r 4
r 2 (9r 2 8) 0 9r 2 8 0 (as r 2 0 is not possible) r 2 2 3 ( ve value of r is not
possible). As d 2V dt 2 0 r 2 2 3 is the point of maxima, i.e., at r 2 2 3 and thus
r2 2 3
Example 2.121 [MN-2007 (2 Solution (d): Speed of dump truck is x km/hr, thus the truck
marks)]: The cost of diesel is Rs. covers one km in 1 x hours; so per km the maintenance cost is
25 ( x 90) per km to drive a Rs. 10 x . So the total cost per km is C 25 x 90 10 x ;
dump truck at a speed of x
km/hour. The maintenance cost of for maximum or minimum, dC dx 0 1 90 10 x 2 0
2 2 3
the truck is Rs. 10 per hour. To x 30 which is the critical point. As d C dx 20 x
minimize the cost per km, the
truck speed in km/hour is d 2C dx 2 20 303 0 x 30 is the point of minima.
x 30
(a) 5 (b) 20 (c) 25 (d) 30
n
Example 2.122 [AG-2008 (2 marks)]: If i1 ( x ai )2 has a minima at A , then A is the arithmetic
mean of the series
(a) a1 a2 a3 ( 1) n 1 a n (b) a1 a2 a3 an
(c) (1 a1 ) (1 a2 ) (1 a3 ) ( 1) n 1 (1 an ) (d) (1 a1 ) (1 a2 ) (1 a3 ) (1 an )
n n
Solution (b): ( x ai ) 2 ( x 2 2ai x ai2 ) nx 2 2( a1 a2 an ) x ( a12 a22 an2 ) .
i 1 i 1
As the given expression on the RHS is a quadratic equation in x . We know that y ax 2 bx c has
a minimum value of D 4a at x b 2a . So if the given expression has a minima at x A then
2( a1 a2 an ) a1 a2 an
A , thus A is the arithmetic mean of a1 a2 a3 an
2n n
Example 2.123 [AG-2008 (1 mark)]: If log e ( y ) x log e ( x ) , then the maximum value of y is
2 1
(a) e (b) e x (c) ee (d) e x
Solution (c): The given function can be written as log e ( y ) log e ( x x ) y x x . For maximum
value of y , y 0 dy dx y (log e x 1) 0 log e x 1 0 x e 1 is critical point for the
given function. As, d 2 y dx 2 y (log e x 1) y x y (log e x 1) 2 1 x x x (log e x 1) 2 1 x
1 1
1 e 1
d 2 y dx 2
x e 1
(e 1 ) e (log e
e 1 1) 2 1 e 1 (e 1 ) e (1 1) 2
e (e ) e0
1
xe is the point of maxima for log e ( y ) x log e ( x ) ; thus maximum value of y is
y x x
x e 1 1 x x 1 (e 1 ) e 1
e e 1
.
Example 2.126 [EC-2008 (1 mark)]: For real values of x , the minimum value of the function
f ( x ) exp( x ) exp( x )
(a) 2 (b) 1 (c) 0.5 (d) 0
x x 2x
Solution (d): For critical values of f ( x ) , f ( x ) 0 e (1 e ) 0 e 1 2 x 0 x 0 .
x
Also, f ( x) e e
x
f ( x ) x 0 e 0 e 0 2 0 x 0 is the point of local minima; thus
minimum value at x 0 is f (0) e 0 e 0 1 1 0 .
[Similar questions were also asked in AE-2007, ME-2007 (1 mark), EC-2014, EE-2014 (2
marks)]
y x 0 16 0 x 0 is the point of minima; y x 2 32 0 x 2 are points of maxima.
Example 2.128 [PI-2008, IN-2007 (2 marks)]: For real x , the maximum value of e sin x e cos x is
(a) 1 (b) e (c) e 2
(d)
sin x cos x sin x cos x 2 sin( x 4)
Solution (c): y e e e e . So the maximum value of y occurs when we
2
have maximum value of sin( x 4) , i.e., 1. Thus maximum value of y is e .
Example 2.129 [MN-2008 (2 marks)]: The volume of a cone is given by V ( 3)l 3 sin 2 cos
where, l is the slant height and is the semi-vertical angle. The angle ( ), for which the volume of
cone becomes maximum is
(a) sin 1 (1 3) (b) cos 1 (1 3) (c) cos 1 ( 2) (d) sin 1 ( 2)
Solution (b): For maximum or minimum, dV d 0 ( 3)l 3 (2 sin cos 2 sin 3 ) 0
( 3)l 3 sin (3 cos2 1) 0 (3cos 2 1) 0 (as sin 0 is not possible) cos 1 3
1 2 2 1
cos (1 3) . As d V d cos1 1 3
0 cos (1 3) is the point of maxima.
1 1 5 1 1 3 1 1 2 1 1
f (t ) t 0 1 2 lim t 3 t t 1 2 lim t
t 0
2! 3! 4! 5! t 0
2! 3! 4! 5!
f (t ) t 0 1 2 (1 2!) (1 3!) 1 (2 3) (1 3) 0 . Hence t 0 is the point of maxima.
Example 2.133 [MT-2010 (2 marks)]: Determine the radius (in m) of a cylinder of volume 200 m3
that has the least surface area
(a) 2.302 (b) 3.142 (c) 3.169 (d) 7.233
2
Solution (c): For a cylinder, Volume V r h ; and whole surface area S 2 r ( r h ) , where r is
the radius of the base and h is its height. As the volume V is given h V ( r 2 )
S 2 r r (V r 2 ) 2 r 2 2V r . For maximum or minimum surface area, dS dr 0
13
4 r 2V r 2 0 r (V 2 )1 3 ; as d 2 S dr 2 13 0 r V (2 ) is the point of
r V (2 )
13
minima as. So r 200 (2 ) 3.169
[Similar question was also asked in TF-2010 (2 marks)]
Example 2.135 [MN-2011 (1 mark)]: The largest area of a rectangular shaft for a given perimeter is
obtained when length is
(a) 2.5 times of breadth (b) 1.5 times of breadth
(c) 2 time of breadth (d) Equal to breadth
Solution (d): Let x and y are the length and breadth of the rectangular shaft; so perimeter
P 2( x y ) and area A xy x ( P 2) x . For maximum or minimum area of the rectangular
shaft dA dx 0 P 2 x x 0 x P 4 . Also d 2 A dx 2 1 ( d 2 A dx 2 ) 1 0
xP 4
Example 2.140 [MT-2013 (2 marks)]: From a 2 m 1.2 m sheet, squares are cut out from each of
the four corners as shown in figure and then the sides are bent to form an open box. The maximum
possible volume (in m3) of the box is ……
Solution: Let box be of length l 2 2 x , breadth b 1.2 2 x and
height h x . So, the volume of the open box is
V (2 2 x)(1.2 2 x ) x 4( x 3 1.6 x 2 0.6 x ) . For maximum and
minimum, dV dx 0 3 x 2 3.2 x 0.6 0 x 0.824, 0.243 . Now
d 2V dx 2 4(6 x 3.2) d 2V dx 2 6.976 0 at x 0.824 ,
x 0.824
2 2
the box has minimum volume; also d V dx 6.968 0 at
x 0.243
x 0.243 , the box has maximum volume, i.e. V 4(0.2433 1.6 0.2432 0.6 0.243) 0.262 m3.
Example 2.141 [EC-2014 (1 mark)]: The maximum value of the function f ( x ) ln(1 x ) x
(where x 1 ) occurs at x …………….
Example 2.142 [EC-2014 (2 marks)]: The maximum value of the determinant among all 2 2 real
symmetric matrices with trace 14 is …………….
Solution: Let x , y be the eigenvalues of a 2 2 real symmetric matrices, then sum of eigenvalues of
any square matrix is equal to its trace x y 14 …(i). Also, product of eigenvalues of a square
matrix gives determinant ( ) of that matrix xy x (14 x ) . So, for maximum or minimum
value of , d dx 0 14 2 x 0 x 7 is the critical point. Also d 2 dx 2 at x 7 is
2 0 , so x 7 is the point of maxima for . Thus 7(14 7) 49 .
Example 2.143 [EC-2014 (2 marks)]: For a right angles triangle, if the sum of the lengths of the
hypotenuse and a side is kept constant, in order to have maximum area of the triangle, the angle
between the hypotenuse and the side is
(a) 12o (b) 36o (c) 60o (d) 45o
Solution (c): From the given triangle y h sin and x h cos ; and as given in the question,
h x k h h cos k h(1 cos ) k , where k constant. so area of the triangle ABC is
given as A (1 2) x y (1 2) h 2 sin cos (1 4) k 2 sin 2 (1 cos ) 2 . For maximum or minimum
area we have to find the critical points of A , i.e.,
dA d 0 (k 2 4) (1 cos ) 2 2 cos 2 (sin 2 )2(1 cos ) sin
(1 cos ) 4 0
(1 cos ) 2 2 cos 2 (sin 2 )2(1 cos ) sin 0
(1 cos ) cos 2 sin 2 sin 0 (as (1 cos ) 0 for the given triangle)
(cos 2 cos 2 cos sin 2 sin ) 0 cos 2 cos(2 ) 0
cos 2 cos 0 2 cos(3 2) cos( 2) 0 cos(3 2) 0 or cos 2 0
cos(3 2) 0 3 2 90o 60o and cos( 2) 0 2 90o 180o
2 2
(which is not possible for the given triangle). Now if we check d A d , we
60o
Example 2.144 [XE-2014 (2 marks)]: The perimeter of a rectangle having the largest area that can
x2 y 2
be inscribed in the ellipse 1 is …………….
8 32
Solution: Let ABCD be a rectangle, of length 2x
and width 2 y , inscribed in a ellipse whose
x2 y2
equation is 1 . As the centre of the
8 32
given ellipse is at the origin, thus any rectangle
inscribe inside the ellipse is symmetrical about
y axis; so coordinates of point A, B, C and D
are shown in figure. So the area of the required
rectangle is A 2 x 2 y 4 xy 4 x 32 4 x 2 ; for maximum or minimum area, we have to find its
32 4 x 2 4 x 2
x (8 x)
critical point, i.e., dA dx 0 4 32 4 x 2
2
0 32 8 x 0
2 2
2 32 4 x 32 4 x
2 2 2
x 2 y 32 4 2 4 . So (2, 4) is the critical points and if we check d A dx , we
(2,4)
will find it is 0 . Thus, (2, 4) is the point of maxima for A and hence perimeter of rectangle ABCD
is P 2(2 x 2 y ) 2(4 8) 24 .
Example 1.146 [AE-2016 (1 mark)]: Let x be a positive real number. The function
f ( x ) x 2 (1 x 2 ) has its minima at x _____.
Solution: f ( x ) 0 2 x ( 2 x 3 ) 0 x 4 1 x 1 . Also f ( x ) 2 (6 x 4 )
f ( x ) x 1 2 6 8 0 x 1 is the point of minima; and f ( x ) x 1 2 6 8 0 x 1
is the point of minima. But x must be positive, so answer is x 1 .
Example 1.149 [EE-2016 (1 mark)]: The maximum value attained by the function
f ( x ) x ( x 1)( x 2) in the interval [1, 2] is _____.
Solution: f ( x) x 2 3x 2 x 2 2 x x 2 x 3 x 2 6 x 2 ;
so f ( x) 0 3 x 2 6 x 2 0 x 1.57, 0.42
So from wavy-curve method, f ( x) changes its sign from ve to ve at x 0.42 .
So x 0.42 is the point of maxima; and x 1.57 is the point of minima. Hence
maximum value of f ( x ) is max{ f (1), f (2)} . As f (1) f (2) 0 . Hence
maximum value of the given function is zero.
[Similar questions were also asked in CS-1997, EC-2007, IN-2008, IN-2012, EE-2012, EC-2012,
EC-2014 (2 marks)]
Example 1.151 [ME-2016 (2 marks)]: Consider the function f ( x ) 2 x 3 3 x 2 in the domain [ 1, 2]
. The global minimum of f ( x ) is _____.
Solution: f ( x ) 2 x 3 3x 2 f ( x ) 6( x 2 x ) f ( x ) 6(2 x 1)
So for critical points of f ( x ) 0 , we have f ( x ) 0 6( x 2 x ) 0 x 0,1
Thus f ( x ) x 0 6 0 x 0 is the point of maxima; f ( x ) x 1 6 0 x 1 is the point of
minima. As our function f ( x) 2 x 3 3x 2 is defined in the closed interval, so global minimum of
f ( x ) min{ f ( 1), f (1), f (2)} min{5, 1, 4} 5 . So answer is ‘ 5 ’.
Example 1.152 [PI-2016 (2 marks)]: The range of values of k for which the function
f ( x ) (k 2 4) x 2 6 x 3 8 x 4 has a local maximum at point x 0 is
(a) k 2 or k 2 (b) k 2 or k 2 (c) 2 k 2 (d) 2 k 2
2 2 3 4
Solution (c): f ( x ) ( k 4) x 6 x 8 x
f ( x ) ( k 2 4)(2 x ) 18 x 2 32 x 3 f ( x ) ( k 2 4)(2) 36 x 96 x 2
Now, at x 0, f ( x) 0 ; so x0 is the point of local maxima if
2 2 2
f ( x ) x 0 0 ( k 4)(2) 36(0) 96(0) 0 k 4 0 k ( 2, 2) .
at least one point P on the curve between A and B at which tangent will be parallel to the x
axis. In Fig. 2.31a, there is only one such point P where tangent is parallel to x axis; in Fig.
2.31b, there are more than one such points where tangents are parallel to the x axis.
Converse of Rolle’s theorem is not true, i.e., if a function f ( x ) is such that f (c) 0 for at least
one c in the open interval ( a, b) , then it is not necessary that (i) f ( x ) is continuous in [ a, b] ; (ii)
3 2
f ( x ) is differentiable in ( a, b) ; and (ii) f (a ) f (b) . For e.g., if f ( x) x x and the interval
is [1,3] , then f ( x ) is continuous in [1,3] and it is differentiable in ( 1,3) but f (1) f (3) .
Lagrange’s Mean Value Theorem: If a function f ( x) is (i) continuous in the closed interval
[ a, b] , and (ii) differentiable in the open interval ( a, b) . Then there is atleast one value c ( a, b) ,
such that; f (c) f (b) f ( a ) (b a )
Geometrical Interpretation: Let f ( x ) be a function defined on [ a, b] and let APB be the curve
represented by y f ( x) . Then co-ordinates of A and B are a, f ( a ) and b, f (b)
respectively. Let the chord AB makes an angle with the axis of x . Then from triangle ARB ,
tan BR AR tan f (b) f (a ) (b a ) . By Lagrange’s Mean value theorem, we have,
f (c ) f (b) f (a ) (b a ) tan f (c ) slope of chord AB slope of tangent at
(c, f (c )) . Thus geometrical meaning of the mean value theorem states: If y f ( x) is continuous
and differentiable in ( a, b) , then at least one point P on the curve in ( a, b) , where the tangent
will be parallel to chord AB . As shown in Fig. 2.32a, there is only one such point P where
tangent is parallel to chord AB ; but in Fig. 2.32b, there are more than one such points where
tangents are parallel to chord AB .
Example 2.154 [ME-1994 (1 mark)]: The value of in the mean value theorem of
f (b) f ( a) (b a ) f ( ) for f ( x ) ax 2 bx c in ( a, b) is
Example 2.155 [CE-2005 (2 marks)]: A rail engine accelerates from its stationary position for 8
seconds and travels a distance of 280 m. According to the Mean Value Theorem, the speedometer at a
certain time during acceleration must read exactly
(a) 0 km/hr (b) 8 km/hr (c) 75 km/hr (d) 126 km/hr
Solution (d): Let the position of rail engine x (t ) is continuous and differentiable function, so
according to Lagranges mean value theorem there exist a ‘ t ’, 0 t 8 such that
x(t ) v (t ) x (8) x (0) (8 0) 280 0 (8 0) 280 8 m/sec (280 8) (18 5) 126 km/hr.
Example 2.159 [XE-2014 (2 marks)]: Let the function f : [0, ) R be such that
2
f ( x ) 8 ( x 3x 4) for x 0 and f (0) 1 . Then f (1) lies in the interval
(a) [0,1] (b) [2, 3] (c) [4, 5] (d) [6, 7]
Solution (b): We have to find the range of f (1) , so let us take the interval x [0,1] , as in x [0,1]
the function is continuous and differentiable so we can apply Lagrange’s MVT, as c [0,1] s.t.
2 2
f (c) { f (1) f (0)} (1 0) f (1) f (0) f (c ) . f (c ) 8 (1 3 1 4) , 8 (0 3 0 4)
(as c [0,1] ). So f (c ) 1, 2 f (1) f (0) 1, f (0) 2 f (1) 2,3 .
[Similar question was also asked in AG-2012 (2 marks)]
Example 2.160 [CH-2016 (2 marks)]: The Lagrange mean-value theorem is satisfied for
f ( x ) x 3 5 , in the interval (1, 4) at a value (rounded off to the second decimal place) of x equal to
___.
Solution: As the given function is continuous in [1, 4] and differentiable in (1, 4) , then there is at
least one value x ( a, b) , such that
f (b) f ( a) 2 f (4) f (1) 2 (43 5) (13 5)
f ( x ) 3x 3x
ba 4 1 4 1
2 2
9 x 63 x 7 x 7 2.64 . So x 2.64 (1, 4) .
Example 2.161 [BT-2013 (1 mark)]: If u log(e x e y ) , then (u x) (u y)
(a) e x e y (b) e x e y (c) 1 (e x e y ) (d) 1
u ex u ey u u
Solution (d):
x
log(e e )
y
x y
and log(e x e y ) x y
1
x x e e y y e e x y
[Similar question was also asked in AG-2014 (1 mark)]
2 f
Example 2.162 [ME-2008 (2 marks)]: Let f y x . What is at x 2 , y 1 ?
xy
(a) 0 (b) ln 2 (c) 1 (d) 1 ln 2
Solution (c):
2 f f y x x 1 x 1 x 1 2 f
( xy ) y xy log e y 12 1 2 12 1 log e 1 1
xy x y x y x xy (2,1)
f 2
( x y 2 z 2 )3 2 3 x 2 ( x 2 y 2 z 2 )1 2
2
( x 2 y 2 z 2 )1 2 ( x 2 y 2 z 2 ) 3x 2
2 2 2 2 3 2 2 2 3
x (x y z ) (x y z )
2 2 2 2
2 f 2 x2 y 2 z 2 f 2y x z 2 f 2z2 x2 y2
. Similarly, and
x
2 2
(x y z )
2 2 52
y 2 ( x 2 y 2 z 2 )5 2 z 2 ( x 2 y 2 z 2 )5 2
2 f 2 f 2 f 2x2 y2 z 2 2 y 2 x2 z2 2z 2 x2 y2
Thus, 0
x 2 y 2 z 2 ( x 2 y 2 z 2 )5 2 ( x 2 y 2 z 2 )5 2 ( x 2 y 2 z 2 )5 2
z z
xy ln( xy ) ln( xy ) ( xy ) xy ln( xy ) ln( xy ) y y x
xy ln( xy ) xy and
x x x x x
z z
xy ln( xy ) ln( xy ) ( xy ) xy ln( xy ) ln( xy ) x x y xy ln( xy ) xy
y y y y y
Example 2.165 [PI-2016 (1 mark)]: For the two functions f ( x, y ) x 3 3xy 2 and
2 3
g ( x, y ) 3 x y y , which one of the following options is correct?
(a) f x g x (b) f x g y (c) f y g x (d) f y g x
f f
Solution (c): ( x3 3 xy 2 ) 3 x 2 3 y 2 , ( x 3 3xy 2 ) 6 xy
x x y y
g g f g
(3 x 2 y y 3 ) 6 xy , (3x 2 y y 3 ) 3x 2 3 y 2 . As , so option (c) is correct.
x x y y y x
z f ( x x, y y ) f ( x, y y ) x f ( x, y y ) f ( x, y ) y
t x t y t
Let t 0 x, y 0 . Also, lim x t dx dt and lim y t dy dt . Hence,
t 0 t 0
dz
lim lim
dt y 0 x 0 f ( x x, y y ) f ( x, y y ) dx
x
dt y
lim
0
f ( x, y y ) f ( x , y ) dy
y
dt
dz f ( x, y y ) dx f ( x, y ) dy f ( x, y ) dx f ( x , y ) dy z dx z dy
lim dt y dt x dt y dt x dt y dt
dt y 0 x
Corollary: If z f ( x, y ) has continuous partial derivatives w.r.t. x and y and x , y both
have continuous partial derivatives w.r.t. u and v , then we have
z z x z y z z x z y
and .
u x u y u v x v y v
2
d y dx dy dx dy 2
f y ( fx ) ( fx ) fx ( f y ) ( f y ) ( f y )
2
dx x dx y dx x dx y dx
d2y
dy dy 2
f y f xx f yx
f x f xy f yy ( fy )
dx
2
dx
dx
d y f f
2
2 f y f xx f yx x f x f xy f yy x ( f y ) 2
dx f y f y
d2y ( f y ) 2 f xx f yx f x f y f x f y f xy f yy ( f x ) 2 ( f y ) 2 f xx 2 f x f y f xy f yy ( f x ) 2
dx 2 ( f y )3 ( f y )3
d 2 y dx 2 ( f y )2 f xx 2 f x f y f xy f yy ( f x )2 ( f y )3 , f y 0 f xy f yx
Total Derivative
If z f ( x, y ) possesses continuous partial derivatives of first order; and x g1 (t ) and y g 2 (t )
, then the total differential of z is given by dz (z x) dx (z y) dy [This point was asked in
ME-2000 (1 mark)].
Proof: Let z f ( x, y ) has continuous partial derivatives w.r.t. x and y ; and x , y have
derivatives w.r.t. t , then we have dz dt (z x)(dx dt ) (z y )(dy dt ) . Multiplying both
sides with dt we get, dz (z x) dx (z y) dy .
In the same way, if z f ( x , y , w) possesses continuous partial derivatives of first order; and
x g1 (t ) , y g 2 (t ) and w g3 (t ) , then the total differential of z is given by
dz (z x) dx (z y) dy (z w)dw .
The above two results can be extended to any number of variables.
If z f ( x, y ) and y g ( x ) then, ( dz dx) (z x) (z y )(dy dx)
Maxima and Minima of Functions in Two Variables: Let z f ( x, y ) be a function with two
independent variables ( x and y ) with continuous second order partial derivatives 2 z x 2 or z xx ,
2 z y 2 or z yy , 2 z (xy ) or z xy and 2 z (yx) or z yx . For a well behaved function
2 z (xy ) 2 z (yx ) or z xy z yx . The critical values, which occur at z x 0 and z y 0 ,
are of three types:
A critical point is a maximum if the value of z at that point is greater than its value at all
sufficiently close surrounding points.
A critical point is a minimum if the value of z at that point is less than its value at all sufficiently
close surrounding points.
A critical point is a saddle point if the value of z at that point is greater than its value at some
sufficiently close surrounding points and less than its value at other sufficiently close surrounding
points.
D 0 and z xx ( x0 , y0 )
0 then ( x0 , y0 ) is the point of maxima.
D 0 and z xx ( x0 , y0 )
0 then ( x0 , y0 ) is the point of minima.
D 0 then z has a saddle point at ( x0 , y0 ) [This point was asked in CE-1998 (1 mark)]
D 0 then no conclusion can be drawn.
2
Example 2.172 [PI-2007 (2 marks)]: The function f ( x, y ) x 2 y 2 defined on R , the point [0, 0]
is
(a) a local minimum (b) neither a local minimum nor a local maximum
(c) a local maximum (d) both a local minimum and a local maximum
Solution (b): As, f x 2 x , f y 2 y , f xx 2 , f yy 2 , f xy 0 . For the critical points of f ( x, y ) ,
we have f x 0 2 x 0 x 0 and f y 0 2 y 0 y 0 ; so (0, 0) is the critical point of
f ( x, y ) . Also D f xx f yy f xy2 4 02 4 0 (0, 0) is a saddle point.
Example 2.173 [XE-2007 (2 marks)]: Let f : R R be twice differentiable real valued function
such that f 1 n 1 for n 1, 2, 3, . Then
(a) f (0) 0 (b) f (0) 1 (c) 0 f (0) 1 (d) f (0) 1
Solution (a): We have f 1 n 1 (d dn) f (1 n) (d dn)1 f (1 n) ( 1 n 2 ) 0 f (1 n) 0 .
So, as n f (1 ) 0 f (0) 0
Example 2.174 [ME-2009 (2 marks)]: The distance between the (a) 3 (b) 3 / 2
origin and the point nearest to it on the surface z 2 1 xy is (c) 1 (d) 2
2
Solution (c): Let P ( x, y , z ) be any point on the surface z 1 xy , then distance between P and
origin is PO ( x 0) 2 ( y 0) 2 ( z 0) 2 x2 y 2 z 2 x 2 y 2 1 xy . So the distance
between origin and point nearest to the surface z 2 1 xy is the minimum value of PO , i.e.,
minimum value of f ( x, y ) x 2 y 2 1 xy . Now f x 2 x y f xx 2 and f y 2 y x f yy 2
, also f xy f yx 1 . So for critical points we have f x 0 2 x y 0 and f y 0 2 y x 0 an
solving these two equation we get x 0, y 0 , so ( x, y ) (0, 0) is the critical point. As the
2 2
discriminant D f xx f yy f xy 4 1 3 0 and f xx (0,0)
2 0 , thus (0, 0) is the point of minima
2 2
for f ( x, y ) and also for PO . Hence PO (0,0) 0 0 1 0 1
Example 2.175 [AE-2010 (2 marks)]: The function (a) (1, 2) (b) (3, 0)
f ( x, y ) x 2 y 2 xy 3 y has an extremum at the point (c) (2, 2) (d) (1, 1)
Solution (a): As, f x 2 x y , f y 2 y x 3 , f xx 2 , f yy 2 , f xy 1 . For critical points we
have f x 0 2 x y 0 and f y 0 2 y x 3 0 ; solving these two equations we have
x 1, y 2 ; so (1, 2) is the critical point. As the discriminant D f xx f yy f xy2 4 1 3 0
(1, 2) is not a saddle point.
1. The distance travelled s (in metre) by a particle in t second is given by s t 3 2t 2 t . The speed
(in cm/s) of the particle after 1 sec will be _____.
2. A man of height 1.8 m is moving away from a lamp post at the rate of 1.2 m/sec. If the height of
the lamp post be 4.5 meter, then the rate (in m/sec) at which the shadow of the man is lengthening
_____.
3. The line x y 2 is tangent to the curve x 2 3 2 y at its point
(a) (1,1) (b) ( 1,1) (c) ( 3, 0) (d) (3, 3)
4. The equation of the normal to the curve y sin( x 2) at (1,1) is
(a) y 1 (b) x 1 (c) y x (d) y 1 (2 )( x 1)
5. If the normal to the curve y f ( x) at the point (3, 4) makes an angle 3 4 with the positive
x axis then f (3) is equal to _____.
6. The point(s) on the curve y 3 3 x 2 12 y where the tangent is vertical (parallel to y axis), is are
(a) ( 4 3 , 2) (b) ( 11 3,1) (c) (0, 0) (d) ( 4 3 , 2)
2 2
7. The angle between the curves y x and x y at (1,1) is
(a) tan 1 (4 3) (b) tan 1 (3 4) (c) 90o (d) 45o
8. The length of sub-tangent to the curve x 2 y 2 a 4 at the point ( a, a ) is
(a) 3a (b) 2a (c) a (d) 4a
9. The sum of intercepts on co-ordinate axes made by tangent to the curve x y a is
(a) a (b) 2a (c) 2 a (d) None of these
10. The length of perpendicular from (0, 0) to the tangent (a) 1 2 (b) 3 5
2
drawn to the curve y 4( x 2) at point (2, 4) is (c) 6 5 (d) 1
5 4 3
11. What are the minimum and maximum values of the function x 5 x 5 x 10 ?
(a) 37, 9 (b) It has 2 minimum and 1 maximum values
(c) 10, 0 (d) It has 2 maximum and 1 minimum values
12. If y a log x bx 2 x has its extremum value at x 1 and x 2 , then ( a, b)
(a) (1,1 2) (b) (1 2 , 2) (c) (2, 1 2) (d) ( 2 3, 1 6)
13. Maximum value of (1 x ) is x
(a) (e) e 1e
(b) (e) (c) (e) e (d) (1 e)e
14. Maximum slope of the curve y x3 3 x 2 9 x 27 is _____.
x
15. The function f ( x ) t (et 1) (t 1) (t 2) 3 (t 3)5 dt has a local minimum at maximum x __.
1
16. If the function f ( x ) 2 x 3 9ax 2 12a 2 x 1 , where a 0 attains its maximum and minimum at
p and q respectively such that p 2 q , then non-zero value of a equals _____.
17. The maximum and minimum values of x 3 18 x 2 96 in interval (0,9) are
(a) 160, 0 (b) 60, 0 (c) 160, 128 (d) 120, 28
2.5 Integration
1 (e
x
Example 2.178 [CH-2013 (1 mark)]: Evaluate 1) dx (Note: c is constant of integration)
(a) e x
(e x 1) c (b) ln (e x
1) e x c (c) ln e x
(e x 1) c x
(d) ln(1 e ) c
dx ex ex 1 ex ex
Solution (d): Let I x
x
dx x 1 dx x dx dx . In first integral,
e 1 e 1 e 1 e 1
x x x
put e x 1 t e x dx dt , I (1 t )dt dx ln t x c ln(e 1) ln e c ln(1 e ) c
Standard Substitution: The following standard substitutions are used according the type of
integrand as:
Integrand Substitution
2 2 2 2 2 2 : x a sin , x a cos
a x , 1 a x , (a x )
: x a tan
x2 a2 , 1 x2 a2 , x2 a 2
: x a sec or x a cosh
x 2 a 2 , 1 x 2 a 2 , x2 a2
x (a x ), ( a x) x , x (a x ),1 x ( a x ) : x a tan 2
x ( x a ), ( x a ) x , x( x a ),1 x( x a) : x a sec
2
(a x) ( a x), (a x ) (a x) : x a cos 2
I I ( a 2 b 2 ) (e ax b 2 )( a sin bx b cos bx ) I e ax ( a 2 b 2 ) (a sin bx b cos bx ) c
evaluate integrals of the above form (i) Make the coefficient of x 2 unity by taking ‘ a ’ common
from ax 2 bx c (ii) Express the terms containing x 2 and x in the form of a perfect square by
adding and subtracting the square of half of the coefficient of x . (iii) Put the linear expression in
x equal to t and express the integrals in terms of t . (iv) The resultant integrand will be either in
standard form which are given as:
dX 1 X A dX 1 X
2 2
ln c 2 2
tan 1 c
X A 2A X A X A A A
px q
Integral of the form 2 dx : Breaking px q into two parts such that one part is the
ax bx c
differential coefficient of the denominator and the other part is a constant. If M and N are two
2
constants, then px q as px q M (d dx)(ax bx c) N M (2ax b) N
px q (2aM ) x Mb N . Comparing the coefficients of x and constant terms on both
sides, we have, p 2 aM M p 2 a and q Mb N N q Mb q ( p 2 a )b . Hence,
px q p 2ax b p dx
the given integral is ax2 bx c dx 2a ax2 bx c dx q 2a b ax 2 bx c , where,
2ax b 2 dx
ax 2 bx c dx log ax bx c ; ax2 bx c is evaluated by the method discussed already.
dx
Integral of the form ax 2 bx c
: To evaluate the given integrals make the coefficient of
x
2
unity by taking a common
from ax 2 bx c . Then,
dx 1 dx b c
. Put x 2 x , by the method of completing
ax 2 bx c a 2
x (b a) x (c a ) a a
2 2 2 2 2 2
the square in the form, A X or X A or X A where, X is a linear function of
x and A is a constant. After this, use any of the following standard formulae according to the
case under consideration
dx x dx
sin 1 c log x x 2 a 2 c
2 2 2 2
a x a x a
dx
log x x 2 a 2 c
2 2
x a
px q
Integral of the form dx : To evaluate this type of integrals, proceed as follows:
ax 2 bx c
d
Let px q M ( ax 2 bx c ) N px q M (2ax b) N , where M and N are
dx
constants. By equating the coefficients of x and constant terms on both sides, we get
p 2 aM M p 2 a and q bM N N q (bp 2 a ) . So, integral breaks up into two
px q p 2ax b bp dx
parts:
ax 2 bx c
dx
ax 2 bx c 2a
dx q
2a ax 2 bx c
I1 I 2 , (say)
p 2ax b 2
Now, I1 dx . Putting ax bx c (2ax b)dx dt , we have,
2a 2
ax bx c
p 1 2 p t1 2 p 2
I1 t dt C1 ax bx c C1 . I 2 is calculated as in the previous
2a 2a (1 2) a
section.
f ( x)
Integrals of the form ax 2 bx c dx , where f ( x ) is a polynomial of degree 2 or greater than
2: To evaluate this type of integrals, divide the Nr by Dr, we get,
f ( x) R( x)
ax 2 bx c dx Q( x)dx ax 2 bx c dx , where, Q( x) is a polynomial and R( x) is a
linear polynomial in x . The integrals on R.H.S. can be obtained by the methods discussed earlier.
x2 1 x2 1
Integrals of the form x4 kx 2 1 dx and x 4 kx 2 1dx : To evaluate the integral of the form
x2 1 1 (1 x 2 )
I dx , divide the Nr and Dr by x 2 to get I dx Put
x 4 kx 2 1 2
x k (1 x )
2
x (1 x) t 1 (1 x 2 ) dx dt and x 2 (1 x 2 ) 2 t 2 x 2 (1 x2 ) t 2 2 . Then, the
dt
given integral reduces to the form I 2
, which can be integrand as usual. To evaluate
t 2k
x2 1 1 (1 x 2 )
I dx, first divide the Nr and Dr by x 2 and we get I dx . Now,
x 4 kx 2 1 2
x k (1 x )
2
2
2 2
let x (1 x) t 1 (1 x ) dx dt and x (1 x ) 2 t x (1 x ) t 2 . Thus, we
2 2 2 2
dt
have t 2
, which can be evaluated as usual.
t 2k
2
Integrals of the forms ax bx cdx : To evaluate this form of integrals, express
2 2 2
ax bx c in the form a ( x ) by the method of completing the square and apply the
standard result discussed in the above section according to the case as may be
(2 ax b ) ax 2 bx c 4ac b 2 dx
ax 2 bx cdx .
4a 8a 2
ax bx c
The following are the result 2 2 2 2 2 2 2
x a dx ( x 2) x a ( a 2) log( x x a ) c
based on the form
2 2 2 2 2 2 1
ax bx cdx : a x dx ( x 2) a x ( a 2) sin ( x a) c
2
Integrals of the form ( px q) ax bx cdx : To evaluate this form of integral, first put,
( px q ) as px q M ( d dx )( ax 2 bx c ) N px q M (2 ax b ) N where, M and N
are constant. Compare the coefficients of x and constant terms on both sides, will get
p 2 a M M p 2 a and q Mb N N q Mb q ( p 2 a )b . Now, the given integral
p pb p pb
ax 2 bx c dx q 2
2a
is, I (2ax b) ax bx c dx I1 q I 2 . To
2a 2a 2a
evaluate I1 , put ax 2 bx c t and to evaluate I 2 , follows the method discussed previously.
dx
Integrals of the form
Q
P
, (where P and Q are linear or quadratic expressions in x ): To
dx
Integral of the form I : To evaluate this type of integrals we substitute
a sin x b cos x
1 dx
a r cos , b r sin and so r a 2 b 2 , tan 1 (b a ) . So, I
r sin( x )
1 1 x 1 x 1 1 b
cosec( x ) dx log tan tan c.
r
I log tan
r 2 2 2
a b 2
2 2 a
1 tan 2 ( x 2) 2 tan( x 2)
The given integral can also be evaluated by using cos x 2
and sin x 2
.
1 tan ( x 2) 1 tan ( x 2)
Integral of the following forms:
dx dx dx dx dx
a b cos2 x , a b sin 2 x , a sin 2 x b cos2 x , (a sin x b cos x)2 , a b sin 2 x c cos 2 x . To
evaluate above forms of integrals divide both the Nr and Dr by cos 2 x . Replace sec 2 x in Dr, if
2 2
any by (1 tan x) . Put tan x t sec xdx dt . Now, evaluate the integral thus obtained, by
the method discussed earlier.
a sin x b cos x
Integrals of the form dx : Such rational functions of sin x and cos x are
c sin x d cos x
integrated by expressing the Nr of the integrand as: Nr M (Differentiation of Dr) N (Dr), i.e.,
d
a sin x b cos x M (c sin x d cos x ) N (c sin x d cos x ) . The arbitrary constants M , N
dx
are determined by comparing the coefficients of sin x and cos x . So, the given integral is
a sin x b cos x M (c cos x d sin x ) N (c sin x d cos x)
I dx dx
c sin x d cos x c sin x d cos x
c cos x d sin x
I M dx N 1dx M log c sin x d cos x Nx c.
c sin x d cos x
a sin x b cos x q
Integrals of the form dx : To evaluate this type of integrals, express the Nr.
c sin x d cos x r
as, (c sin x b cos x q) M (c sin x d cos x r ) N (c cos x d sin x ) P , where M , N , P are
constants to be determined by comparing the coefficients of sin x, cos x and constant term.
a sin x b cos x q Diff.of denominator dx
dx M dx N dx
c sin x d cos x r Denominator c sin x d cos x r
a sin x b cos x q dx
dx Mx N log Denominator P
c sin x d cos x r c sin x d cos x r
m n
Integral of the form sin x cos x dx : (i) To evaluate the integrals of the form
polynomial ( x 5) and a proper function (19 x 30) ( x 2 5 x 6) .
Partial fractions: Any proper rational function can be broken up into a group of different rational
fractions, each having a simple factor of the denominator of the original rational function. Each such
fraction is called a partial fraction. If by some process, we can break a given rational function
f ( x) g ( x) into different fractions, whose denominators are the factors of g ( x ) , then the process of
obtaining them is called the resolution or decomposition of f ( x) g ( x) into its partial fractions.
Depending on the nature of the factors of the denominator, the following cases arise.
When the denominator consists of non-repeated linear factors: To each linear factor ( x a )
occurring once in the denominator of a proper fraction, there corresponds a single partial fraction
of the form A ( x a ) , where A is a constant to be determined by equating the coefficients of like
powers of x in the resulting identity and solving the equations.
When the denominator consists of linear factors, some repeated: Each linear factor ( x a )
occurring r times in the denominator of a proper rational function, corresponds a sum of r
A1 A2 Ar
partial fractions of the form, 2
r
, where A ' s are constants to be
( x a ) ( x a) (x a)
determined by equating the coefficients of like powers of x in the resulting identity and solving
the equations.
When the denominator consists of quadratic factors: To each irreducible non repeated
2 Ax B
quadratic factor ax bx c, there corresponds a partial fraction of the form 2 , where
ax bx c
A and B are constants to be determined. To each irreducible quadratic factor ax 2 bx c
occurring r times in the denominator of a proper rational fraction there corresponds a sum of r
A1 x B1 A2 x B2 Ar x Br
partial fractions of the form 2
2 2
2 r
, Where,
( ax bx c ) ( ax bx c ) ( ax bx c )
A1 , A2 ,, Ar and B1 , B2 , Br are constants to be determined by equating the coefficients of like
powers of x in the resulting identity and solving the equations.
3x 1
Example 2.182: Evaluate 2
( x 2) ( x 2)
Solution:
3x 1 A B C
2
2
3 x 1 A( x 2) ( x 2) B( x 2) C ( x 2) 2 …(i);
( x 2) ( x 2) ( x 2) ( x 2) ( x 2)
Putting x 2, 2 successively in equation (i), we get B 7 4 , C 5 16 . Now, we put x 0 and get
3x 1 5 dx 7 dx 5 dx
A 5 16 . I 2
dx 2
( x 2) ( x 2) 16 x 2 4 ( x 2) 16 x 2
5 7 1 5 5 x2 7
I log( x 2) log( x 2) c log c
16 4 ( x 2) 16 16 x2 4( x 2)
2x2 3 x 1 1 x c
Example 2.183: If ( x2 1)( x 2 4) dx a log x 1 b tan
2
then find values of a and b .
2 2x2 3 2y 3 2y 3 A B
Solution: Put x y 2 2
( x 1)( x 4) ( y 1) ( y 4) ( y 1) ( y 4) ( y 1) ( y 4)
2 y 3 A( y 4) B ( y 1) . Comparing the coefficient of y and constant terms A B 2 ,
1 1 1 1
4 A B 3 A 1 , B 1 . I dx dx 2
dx 2
dx
y 1 y4 x 1 x 4
1 x 1 1 x
I log tan 1 c a 1 2 and b 1 2 .
2 x 1 2 2
Some integrals which cannot be found: Any function continuous on interval (a, b) has an anti-
derivative in that interval. In other words, there exists a function F(x ) such that F' ( x ) f ( x ).
However not every anti-derivative F ( x ), even when it exists is expressible in closed form in terms of
elementary functions such as polynomials, trigonometric, logarithmic, exponential etc. function. Then
we say that such anti-derivatives or integrals ‘cannot be found.’ Some typical examples are:
dx x2 cos x
3 2
1 x dx e dx sin( x 2 ) dx dx x tan x dx
log x x
sin x x2 x2 2
x dx 1 x5 dx e dx cos( x ) dx sin x dx 1 x3 dx
b
If a f ( x ) dx 0 , then the equation f ( x ) 0 has at least one root lying in ( a, b) , provided f is
a continuous function in ( a, b) .
Evaluation of Definite Integral by Substitution Method: When the variable in a definite
integral is changed, the substitutions in terms of new variable should be effected at three places:
(i) In the integrand, (ii) In the differential say, dx, (iii) In the limits. For example, if we put
b b g (b )
g ( x) t in the integral a f { g ( x )}g ( x ) dx , then a f { g ( x )} g ( x ) dx g ( a ) f (t ) dt .
1
Example 2.184 [CE-2002 (2 marks)]: Evaluate 0 x ln x dx .
(a) 1 4 (b) 0 (c) 1 4 (d) 1
Solution (c): By using product rule of integration, let us take ‘ ln x ’ as 1 function and ‘ x ’ as 2nd
st
2 2
d ln x
function; thus I x ln
x dx ln x xdx
dx
xdx dx (ln x) x2 x4 c
2nd 1st
1
1 x2 x2 1 1 0 0 1
Thus, x ln xdx (ln x) c (ln 1) c (ln 0) c
0 2 4 2 4 2 4 4
0
[Similar question was also asked in EE-2010 (1 mark)]
3
Example 2.185 [EC-2006 (2 marks)]: 0 sin d is given by
(a) 1 2 (b) 2 3 (c) 4 3 (d) 8 3
3 3
Solution (c): As sin 3 3sin 4sin sin (3sin sin 3 ) 4 , thus
3 1 cos 3 1
1 1 1 4
0 sin d 4 0 (3sin sin 3 )d 4 3cos 3 0 4 3 3 3 3 3
[Similar question was also asked in CE-2001 (1 mark)]
1
I ( x 2) ( x 2)
2
(
0
2) (1 2) 0 0 ( 2) (1 2)
Example 2.189 [CH-2010 (2 marks)]: For a function g ( x ) , if g (0) 0 and g (0) 2 , then
g ( x)
lim (2t x ) dt is equal to
x 0 0
2
I sin 0 ( 2 1) (0 0) ( 2) 1
5
Example 2.191 [AE-2013 (1 mark)]: The value of 4 {( x 2) ( x 2 4 x 21)}dx is
2 x3 2 2 x1 2 2 x3 2 2 2 x3 2
I (ln x ) dx (ln x ) c . Thus,
3 3 3 3 3
e
e 2 x3 2 4 x3 2 2e 3 2 4 e3 2 2 4 2e 3 2 4
1 x ln x dx (ln x )
3
9
c (ln e)
1 3
9
c (ln 1)
3
9
c
9
9
2
2 cos x
Example 2.194 [AG-2014 (2 marks)]: The value of 0
1 sin x
dx is
d
Solution: x sin x dx x sin xdx dx x sin xdx x cos x cos xdx x cos x sin x c
2
So 0 x sin x dx ( x cos x sin x c)0 2 (0 1 c) (0 0 c) 1
e
Example 2.198 [PE-2016 (1 mark)]: The value of the definite integral 1 (ln x)dx is equal to _____.
d
Solution: 1(ln x ) dx ln x 1dx
dx
(ln x ) 1dx x ln x 1dx x ln x x
e e
So 1 (ln x)dx ( x ln x x)1 (e ln e e) (1ln1 1) 0 (0 1) 1 .
4 x2 1
Example 2.199 [AG-2016 (2 marks)]: The value of the integral, I dx is _____.
2 x2 1
x2 1 x2 1 2
Solution: I
2
4
2
x 1
dx
2
4
x 12
4
dx 1dx
2 2
4
x 12
2
dx ( x ) 42
4
2 1
x 1
1
x 1 dx
4
I (4 2) log( x 1) log( x 1)2 2 log 3 log 5 log1 log 3 2 log(9 5) 2.25
Example 2.200 [AE-2016 (2 marks)]: The value of definite integral 0 ( x sin x)dx is _____.
Solution: I ( x sin x) dx {( x) sin( x)}dx {( x) sin x}dx
0 0 0
I sin x dx x sin x dx ( cos x)0 I 2 I 2 I 3.14
0 0
(c)
(d)
b
Geometrical Interpretation of Definite Integral: In general, a f ( x ) dx represents an algebraic
sum of areas of the region bounded by the
curve y f ( x) , the x axis and the ordinates
x a and x b , as shown in Fig. 2.33a. The
area above the x axis are taken positive,
while those below the x axis are taken
negative. For e.g., in Fig. 2.33b,
b
a f ( x ) dx A1 A2 A3 .
b
Similarly, a f ( y ) dy represents an algebraic
sum of areas of the region bounded by the
curve x f ( y ) , the y axis and abscissa
y a and y b , as shown in Fig. 2.33c. The
area left the y axis are taken positive, while
those right the y axis are taken negative. For
b
Figure 2.33: Geometrical interpretation of Definite integral
e.g., in Fig. 2.33d, a f ( y ) dy A1 A2 A3
Definite Integral as the Limit of Sum: Let f be a continuous function (whose graph lies in 1st
b
quadrant) defined on a closed interval [ a, b] . The definite integral a f ( x ) dx is the area bounded by
the curve y f ( x) , the ordinates x a , x b and the x axis. To evaluate this area, consider the
region FPDEF between the curve, x axis and the ordinates x a , x b , as shown in Fig. 2.34.
Divide the interval [ a, b] into n equal sub-intervals denoted by [ x0 , x1 ],[ x1 , x2 ],,
[ xr 1 , xr ],,[ xn1 , xn ] , where, x0 a, x1 a h, x2 a 2h,, xr a rh,, xn a nh b
n (b a ) h , where, h is separation between the successive intervals. From the figure, area of
Some useful results for evaluation of definite integrals as limit for sums
1 2 3 n n ( n 1) 2 2 2 2 2
1 2 3 n n(n 1)(2n 1) 6
2
13 23 33 n 3 n (n 1) 2 a ar ar n1 a(r n 1) (r 1) , r 1
cos (ei ei ) 2 , sin (e i e i ) 2 cosh (e e ) 2 , sinh (e e ) 2
1 1 1 1 1 2 1 1 1 2
1
22 32
52 62 6 42 24 22 42 62
1 2
1 1 1 1 1 1 2
1 2 2 1 2 2 2 2 2
3 5 8 2 3 4 5 6 12
n 1
n 1 nh h
sin a sin[ a (n 1)h] [sin(a nh)] sin a h sin sin
r 0 2 2 2
n 1
n 1 h sin nh sin h
cos a cos[ a (n 1)h] [cos(a nh)] cos a
r 0 2 2 2
2
Example 2.202: If S n 1 (1 n ) 1 (2 2n ) 1 ( n n ) then find lim Sn .
n
Solution:
1 1 1 1
lim Sn lim lim dx 2[log(1 x )]10 2 log 2
n n r rn n
n (r n) r n 0
x (1 x )
n 1 n x [ x] (a) ln 10 1
Example 2.203 [MA-2016 (2 marks)]: Let S n k 1 and I n 2
dx . (b) ln10 1
k 1 x
Then S10 I10 is equal to (c) ln10 (1 10)
(d) ln10 (1 10)
10 x [ x] 10 1 10 [ x ] 10 [ x ] 10 [ x ]
Solution (a): I10 2
dx dx 2
dx (ln x)110 2
dx ln10 dx …(i)
1 x 1 x 1 x 1 x 1 x2
10 [ x ] 2 [ x] 3 [ x] 4 [ x] 5 [ x] 10 [ x]
Now let 1 2
dx 2
dx 2
dx 2
dx 2
dx dx
x 1 x 2 x 3 x 4 x 9 x2
10 [ x ] 2 1 3 2 4 3 5 4 10 9
2
dx 2
dx 2
dx 2
dx 2
dx dx
1 x 1 x 2 x 3 x 4 x 9 x2
2 3 4 5 10
10 [ x] 1 2 3 4 9
dx
1 x2
x 1 x 2 x 3 x 4 x 9
10 [ x ] 1 2 3 4 9 1 1 1 1 1
dx 1 1 1 1 1
1 x2
2 3 4 5 10 2 3 4 5 10
10 [ x ] 10 1
dx k 1 1 S10 1 . Thus from (i), I10 ln10 ( S10 1) S10 I10 ln10 1 .
2
1 x k
Properties of Definite Integral
b b
Property 1: Changing Dummy Variables a f ( x ) dx f (t ) dt , i.e. the value of a definite
a
integral remains unchanged if its variable is replaced by any other symbol provided the limits of
integration remains the same.
b a
Property 2: Interchanging Limits a f ( x )dx f ( x ) dx , i.e. by the interchange in the limits
b
of definite integral, the sign of the integral is changed.
b c b
Property 3: Splitting Limits a f ( x) dx f ( x) dx f ( x ) dx , i.e. where c may lies inside or
a c
b
outside the interval [ a, b] . Proof: Let f ( x)dx F ( x) a f ( x ) dx F (b ) F ( a ) and,
c b
a f ( x ) dx f ( x ) dx F (c ) F ( a ) F (c) F ( a ) F (b) F (a ) . Here we have not used any
c
condition whether c ( a, b) or c ( a, b) . This property is useful when f ( x ) is not continuous in
[ a, b] because we can break up the integral into several integrals at the points of discontinuity so
that the function is continuous in the sub-intervals.
1 dx
Example 2.204 [AE-2010 (1 mark)]: The definite integral 2
1 x
(a) does not exist (b) is equal to 2 (c) is equal to 0 (d) is equal to –2
2
Solution (d): 1 x is not continuous in ( 1,1) , but continuous in (1,0) and (0,1) . Thus,
0 1
1 dx 0 dx 1 dx 1 1 1 0 1
1 1 1 1 1
1 x 2
1 2
0 2
2 .
x x x 1 x 0 x 1 x 0 0 1 1 0
2
Example 2.205 [CS-2014 (1 mark)]: If 0 x sin x dx k , then the value of k is equal to ………
2 2
x sin x, x [0, ]
Solution: x sin x x sin x x sin x dx x sin x dx . Using product
x sin x, x [ , 2 ] 0 0
dx
rule, x sin x dx x sin x dx sin x dx dx x cos x cos x dx x cos x sin x c
dx
2
x sin x ( x cos x sin x c )0 ( x cos x sin x c )2 4 k k 4 .
0
Example 2.206 [CS-2014 (2 marks)]: The value of the integral 0 x 2 cos x dx is
(a) 2 (b) (c) (d) 2
Solution (a): Using product rule of integration,
d
x 2 cos x dx x 2 cos x dx x 2 cos x dx dx x 2 sin x 2 x sin x dx
dx
d
x sin x dx x sin x dx dx x sin x dx dx x cos x cos x dx x cos x sin x
4 4 4 1 1 1
I tan 0 x dx tan x dx ln cos x 0
ln ln(1) ln ln 2
0 0
4 4 2 2 2
3 10 sin x
Example 2.209 [MN-2010 (2 marks)]: The value of the given integral 5 sin x cos x
dx is
3 10 sin{(3 10) ( 5) x}
I dx
5 sin{(3 10) ( 5) x} cos{(3 10) ( 5) x}
3 10 sin( 2 x ) 3 10 cos x
I dx dx
5 sin( 2 x ) cos( 2 x ) 5 cos x sin x
3 10 sin x 3 10 cos x 3 10 sin x cos x 3 10 3
2I dx dx dx dx
5 sin x cos x 5 cos x sin x 5 cos x sin x 5 10 5 10
I 20
a x
Example 2.210 [CE-2011 (2 marks)]: What is the value of the definite integral, 0 x ax
dx ?
2 (1 x ) 2 sin(1 x ) 2 ( x 1) 2 sin( x 1)
I dx dx I I 2 I 0 I 0
0 (1 x ) 2 cos(1 x ) 0 ( x 1) 2 cos( x 1)
a a
Property 5: a f ( x ) dx { f ( x ) f ( x)}dx , this property is generally used when integrand is
0
either even or odd function of x .
a 0 a
Proof: a f ( x)dx a f ( x)dx 0 f ( x ) dx . Put, x t in first term on R.H.S. dx dt ;
a 0 a
also, x a t a; x 0 t 0 f ( x) dx f ( t )( dt ) f ( x ) dx
a a 0
a a a a a
f ( x )dx f ( t ) dt f ( x) dx f ( x ) dx f ( x ) dx
a 0 0 0 0
a
2 f ( x ) dx, if f ( x ) is an even function or f ( x ) f ( x)
a
a
f ( x ) dx 0
0
, if f ( x) is an odd is odd function or f ( x ) f ( x )
[This property was asked in ME-2011 (1 mark)]
x
If f ( t ) is an odd function, then g ( x ) f ( t ) dt is an even function.
a
x x
Proof: g ( x ) f (t ) dt g ( x ) f (t ) dt
a a
x x
Let t y g ( x ) f ( y )( dy ) f ( y )dy [ f ( y ) is an odd function]
a a
a x x
g ( x) f ( y )dy f ( y )dy 0 f ( y ) dy g ( x ) g ( x) is an even function.
a a a
x
If f ( t ) is an even function, then for a non-zero ‘ a ’, a f ( t )dt is not necessarily an odd
a
function. It will be odd function if 0 f ( t )dt 0 .
x x
Proof: g ( x ) f (t ) dt g ( x ) f (t ) dt
a a
x x
Let t y g ( x ) f ( y )( dy ) f ( y ) dy [ f ( y ) is an even function]
a a
a x 0 a x
g ( x) f ( y ) dy f ( y )dy f ( y ) dy f ( y )dy f ( y )dy
a a a 0 a
a 0 a
Now if 0 f (t ) dt 0
a
f ( y )dy 0, f ( y ) dy 0
0
x
g ( x) f ( y ) dy g ( x ) . Hence, g ( x ) is an odd function.
a
x
If f ( t ) is an even function, then g ( x ) a f ( t )dt is an odd function. Proof is similar to
the above property.
Example 2.212 [CS-1998 (5 marks)]: Evaluate x cos x dx
Solution: Let f ( x ) x cos x f ( x) ( x ) cos( x ) x cos x f ( x) f ( x ) is an odd function
2
and thus 2 f ( x)dx 0 .
2 sin 2 x
Example 2.213 [CE-2002 (1 mark)]: Evaluate 2 1 cos x dx
2
(a) 2 ln 2 (b) 2 (c) 0 (d) (ln 2)
sin 2 x sin( 2 x ) sin 2 x
Solution (c): Let f ( x) f ( x) f ( x ) f ( x ) is an odd
1 cos x 1 cos( x ) 1 cos x
2
function and thus 2 f ( x)dx 0 .
[Similar question was also asked in PI-2008 (1 mark)]
a 6
Example 2.214 [ME-2005 (1 mark)]: a (sin x sin 7 x ) dx
a a a
(a) 2 sin 6 x dx (b) 2 sin 7 x dx (c) 2 (sin 6 x sin 7 x ) dx (d) zero
0 0 0
6 7
Solution (a): Let f ( x) sin x sin x , then using property 5 of definite integral, we have
a 6 a
a (sin x sin 7 x ) dx {sin 6 x sin 7 x sin 6 ( x ) sin 7 ( x)}dx
0
a a
(sin 6 x sin 7 x sin 6 x sin 7 x) dx 2 sin 6 x dx
0 0
5 2 2
Example 2.215 [XE-2011 (1 mark)]: The integral 5 2 f ( x)dx , where f ( x ) e x sin 3 x 4 cos x ,
equals
(a) 4 (b) 8 (c) 5 2 (d) 5 2
5 2 x 2 5 2 2 5 2
Solution (a): 5 2 (e sin 3 x 4 cos x) dx (e x sin 3 x ) dx (4 cos x )dx
5 2 5 2
5 2 2 5 2 2 2
Let I1 (e x sin 3 x ) dx and I2 (4 cos x ) dx . As e ( x ) sin 3 ( x ) e x sin 3 x
5 2 5 2
x2 3
e sin x is an odd function I1 0 ; also 4 cos( x) 4 cos x 4 cos x is an even function
5 2
I2 2 (4 cos x )dx (4 sin x)50 2
4 sin(5 2) sin(0) 4(1 0) 4 I 4
0
2a a a
Property 6: 0 f ( x ) dx f ( x) dx f (2a x )dx
0 0
2a 0 , if f (2a x ) f ( x )
In particular, 0 f ( x ) dx a
2 0 f ( x ) dx , if f (2a x ) f ( x )
It is generally used to make half the upper limit.
2a a 2a
Proof: 0 f ( x) dx f ( x ) dx f ( x) dx , putting x 2a t dx dt in the second
0 a
integral on R.H.S. Also, x a t a and x 2a t 0 .
2a a 0 a a
f ( x ) dx f ( x) dx f (2a t )( dt ) f ( x) dx f (2a t ) dt
0 0 a 0 0
2a a a
f ( x ) dx f ( x) dx f (2a x) dx
0 0 0
2
Example 2.216 [MN-2013 (2 marks)]: The value of 0 log(cos x ) dx is
(a) ( 2) log 2 (b) ( 4) log 2 (c) ( 2) log 2 (d) ( 4) log 2
2 2
f ( x )dx b f (a b x ) dx
b
a a
Solution (a): Let I log cos xdx I log sin xdx
0 0
2 2 2
2I
0
(log sin x log cos x )dx
0
log(sin x cos x )dx
0
log (sin 2 x) 2 dx
2 2 1
2I log(sin 2 x ) dx [Putting 2 x t dx dt 2 ]
2 0
log 2dx log(sin t ) dt log 2
0 0 2
1 2 2 a f ( x ) dx 2 a f ( x )dx, if f (2a x ) f ( x )
0 0
2I 2 log(sin t ) dt log 2
2 0 2
2 b f ( x ) dx b f (t ) dt
a a
2I log(sin x) dx
log 2
0 2
1
2 I I log 2 I log 2 log (log1 log 2) log 2
2 2 2 2 2 2
2a a
Property 7: 0 f ( x )dx { f ( a x ) f (a x )}dx
0
a a a
Proof: R.H.S. { f ( a x ) f ( a x )}dx f ( a x ) dx f ( a x ) dx ,
0 0 0
a aa
R.H.S. f a ( a x ) dx
0 0 a
f ( x )dx In second integral replace x a by x
a 2a 2a
R.H.S. f ( x ) dx f ( x )dx f ( x ) dx L.H.S.
0 a 0
b 1
Property 8: a f ( x )dx ( b a ) f {( b a ) x a )}dx
0
1
Proof: R.H.S. (b a ) f (b a ) x a dx
0
nT T
Property 9: If f ( x ) is a periodic function with period T , then 0 f ( x )dx n f ( x )dx ,
0
where n I .
nT T 2T nT
Proof: 0 f ( x ) dx 0 f ( x ) dx f ( x ) dx f ( x ) dx
T ( n 1) T
nT T T T
f ( x )dx f ( x ) dx f ( x T )dx f x (n 1)T dx
0 0 0 0
nT T T
f ( x )dx f ( x ) dx f ( x ) dx n (times)
0 0 0
nT T
f ( x ) f ( x T ) f x ( n 1)T
0
f ( x )dx n f ( x) dx
0
0 nT a nT T
I f ( x) dx f ( x ) dx f ( x ) dx f ( x) dx n f ( x ) dx
a 0 0 0 0
a nT
Hence, if f ( x ) is a periodic function with period T , then, we can notice that a f ( x) dx is
independent of a .
Property 11: If f ( x) is a periodic function with period T, then
nT T
mT f ( x )dx ( n m )0 f ( x )dx , n,m I
Proof: [Putting x y mT dx dy , x mT y 0; x nT y ( n m)T ]
nT ( n m )T ( n m )T
Hence, L.H.S. f ( x ) dx f ( y mT ) dy f ( y ) dy
mT 0 0
T T
( n m) f ( y ) dy ( n m) f ( x ) dx R.H.S. f ( y mT ) f ( y )
0 0
b nT b
Property 12: If f ( x ) is a periodic function with period T , then a nT f ( x )dx f ( x )dx ,
a
n I
Proof: [Putting x y nT dx dy , x a nT y a; x b nT y b ]
b nT b b b
L.H.S. f ( x )dx f ( y nT )dy f ( y )dy f ( x ) dx R.H.S. f ( y nT ) f ( y )
a nT a a a
Gamma Function:
2 (m 1) 2 ( n 1) 2
If m, n non-negative integers, then 0 sin m x cos n xdx
, where, (n) is
2 ( m n 2) 2
called gamma function which satisfied the following properties: (i) ( n 1) n( n) n ! ; (ii)
(1) 1 ; (iii) (1 / 2) .
In place of gamma function, we can also use the following formula :
2 m n (m 1)( m 3) (2 or 1)(n 1)(n 3) (2 or 1)
0 sin x cos xdx ( m n)( m n 2) (2 or 1)
Note that, it is important to note that we multiply by ( 2) ; when both m and n are even.
d v( x)
Proof: Let F ( x) f ( x) f (t ) dt F {v ( x)} F {u ( x)}
dx u ( x)
d v( x ) d d d
f (t ) dt F {v ( x )} F {u ( x)} F {v ( x )} v ( x) F {u ( x )} u ( x)
dx u ( x ) dx dx dx
d v(x) d d
f (t )dt f {v ( x )} v ( x) f {u ( x )} u ( x )
dx u ( x ) dx dx
If functions ( x) and ( x) are defined on [a, b] and differentiable at a point x ( a, b), and
f ( x, t ) is continuous, then,
d ( x)
dx ( x )
f ( x , t ) dt
( x ) dx dx
( x ) d f ( x, t ) dt d ( x) f ( x, ( x)) d ( x) f ( x, ( x))
dx
Let f ( x, t ) be a function s.t. both f ( x, t ) and its partial derivatives f x ( x, t ) are continuous in t
and x in some region of ( x, t ) plane, including a( x) t b( x ) , x0 x x1 . Let the functions
a ( x ) and b ( x ) are both continuous and both have continuous derivatives for x0 x x1 . Then
d
for x0 x x1 ,
dx
b( x )
a(x)
f ( x, t ) dt f x, b ( x ) b( x ) f x , a ( x ) a ( x )
b( x )
a(x)
f x ( x, t ) dt .
x2
Example 2.217 [XE-2007 (2 marks)]: Let f ( x ) sin tdt for x 0 . Then f ( / 2) is equal to
0
Important results
Every continuous function defined on [a, b] is integrable over [a, b] .
Every monotonic function defined on [a, b] is integrable over [a, b] .
b
If f ( x ) is a continuous function on [a, b] , then c ( a, b) such that a f ( x) dx f (c )(b a ) .
b
f (c ) 1 (b a) f ( x) dx is called the mean value of f ( x ) for the interval [a, b] .
a
If at every point x of an interval [a, b] , the inequalities g ( x ) f ( x ) h( x ) satisfies, then
b b b
a g ( x)dx a f ( x) dx h ( x) dx , where, a b .
a
If m and M are the smallest and greatest values of a function f ( x ) on an interval [a, b] , then
b
m(b a ) f ( x )dx M (b a ) .
a
b b b b
Proof: m f ( x ) M mdx f ( x ) dx Mdx m (b a ) f ( x ) dx M (b a )
a a a a
b b
a f ( x)dx f ( x) dx
a
b b b
Proof: f ( x) f ( x) f ( x) , x [a, b] f ( x ) dx f ( x )dx f ( x) dx
a a a
b b b b b
f ( x ) dx f ( x ) dx f ( x ) dx a f ( x)dx f ( x) dx
a a a a
a
b
f ( x ) g ( x ) dx a
b
f 2 ( x ) dx g ( x) dx , if f
b
a
2 2
( x) and g 2 ( x) are integrable on [a, b]
2 b
Proof: Let, F ( x ) f ( x ) g ( x) 0, R
a
f ( x) g ( x)2 dx 0
b b b
2 g 2 ( x ) dx 2 f ( x ) g ( x) dx f 2 ( x )dx 0 , hence, its discriminant is non-positive.
a a a
2
4 a
b
4 f ( x) dx g ( x) dx
f ( x) g ( x ) dx
b
a
2 b
a
2
f ( x ) g ( x ) dx f ( x) dx g ( x ) dx
b b 2 b 2
a a a
1
2 2
1 x 10 . Then f ( x) x dx 1 x dx satisfies
2 2 2
does not exist on the left side of y axis. Similarly the curve a y x (a x) does not exist for
x a as the values of y are imaginary for x a.
2 2
Example 2.219 [CE-1997 (1 mark)]: The curve given by the equation x y 2axy , is
(a) symmetric about x axis (b) symmetric about y axis
(c) symmetric about the line y x (d) tangential to x y a 3
Solution (c): As by interchanging x and y , the given equation remains unchanged, so the curve is
symmetrical about the line y x .
Solution (a): For the given function, there is no symmetry along any axis or any line. As x ,
F ( x) 2 0 thus option (c) and option (d) are not correct; also at x 0 , F ( x ) 0 , which is
2 2
1(k1 x k 2 x 1) x(2k1 x k2 ) k1 x 1
satisfied by options (a) and (b). Now, F ( x) 2 2
2 2
( k1 x k2 x 1) ( k1 x k 2 x 1)
F ( x) 0 if k1 x 2 1 0 x 1 k1 1 k1 . So only option (a) has zero slope at x 1 k1 ,
hence option (a) is correct and option (b) is wrong.
The length of arc of the curve in parametric form, i.e., x f (t ) and y f (t ) from t t1 to t t2
t2
is given as s ( dx dt ) 2 (dy dt ) 2 dt
t1
The length of arc of the curve in polar form, i.e., r f ( ) and from 1 to 2 is given as
1
s r 2 ( dr d ) 2 d
2
32
Example 2.221 [ME-2008 (2 marks)]: Length of the curve y (2 3) x between x 0 and x 1 is
(a) 0.27 (b) 0.67 (c) 1 (d) 1.22
dy d 2 3 2 b 1 1
Solution (d): x x1 2 s 1 ( dy dx ) 2 dx 1 ( x1 2 ) 2 dx 1 xdx
dx dx 3 a 0 0
1
Now, let I (1 x ) dx , putting 1 x t dx dt and at x 0 , t 1 and at x 1 , t 2 ; thus,
0
2 2
I
1
tdt (2 3) t 3 2
1
(2 3) (2)3 2 (1)3 2 1.219 . Hence s 1.219 1.22
Example 2.222 [CE-2010 (2 marks)]: A parabolic cable is held between two supports at the same
level. The horizontal span between the support is L . The slag at the mid-span is h . The equation of
2 2
the parabola is y 4h( x L ) , where x is the horizontal coordinate and y is the vertical coordinate
with the origin at the centre of the cable. The expression for the total length of the cable is
L 2 2 4 L
(a) 0 1 64h x (1 L )dx (b) 2 1 64 h3 x 2 (1 L4 ) dx
0
L2 2 2 4 L2
(c) 0 1 64h x (1 L )dx (d) 2 1 64 h 2 x 2 (1 L4 ) dx
0
Solution (d): As the length of the curve y f ( x) between x a and x b is given by
b 2 2 2
L 1 ( dy dx ) 2 dx . So y 4h( x L ) dy dx 8h( x L ) ; also as the origin is at the centre of
a
Example 2.223 [AG-2009 (1 mark)]: A curve is having the equation r a(1 cos ) . The perimeter
of the curve between 0 to 2 is
(a) 2a (b) 4a (c) 6a (d) 8a
2
Solution (d): dr d {d d }a (1 cos ) a sin s r 2 ( dr d ) 2 d
1
2 2 2
s a 2 (1 cos ) 2 (a sin ) 2 d 2 a sin( 2) d (2 a)2 cos( 2) 0 4 a(1 1) 8a
0 0
Sign convention for finding the Areas using Integration: While applying the discussed sign
convention, we will discuss the three cases.
b
Case I: In the expression a f ( x ) dx if b a and f ( x ) 0 for all a x b , then this
integration will give the area enclosed between the curve f ( x ) , x axis and the line x a
and x b which is ve . No need of any modification.
b
Case II: In the expression a f ( x ) dx if b a and f ( x ) 0 for all a x b , then this
integration will calculate to be ve . But the numerical or the absolute value is to be taken to
mean the area enclosed between the curve y f ( x) , x axis and the lines x a and x b .
b
Case III. If in the expression a f ( x ) dx where b a but f ( x ) changes its sign a numbers
of times in the interval a x b , then we divide the region [ a, b] s.t. we get the points lying
between [ a, b] where f ( x ) changes its sign. For the region where f ( x ) 0 we integrate to
get the area in that region and then add the absolute value of the integration calculated in the
region where f ( x ) 0 to get the desired area between the curve y f ( x) , x axis and the
line x a and x b . Hence as shown in Fig. 2.37, the shaded area is given by,
c d e f b
A f ( x) dx c f ( x ) dx f ( x ) dx e f ( x ) dx f ( x ) dx
a d f
Symmetrical Area: If the curve is symmetrical about a coordinate axis (or a line or origin), then
we find the area of one symmetrical portion and multiply it by the number of symmetrical
portions to get the required area.
2 2 2
Example 2.224: Find the whole area of circle x y a .
Solution: The required area is symmetric about both the axis as shown in figure
a a
∴ Required area a 2 x 2 dx 4 ( x 2) a 2 x 2 (a 2 2) sin 1 ( x a) 4( 2)(a 2 2) a 2 .
0 0
Area between Two Curves
When both curves intersect at two points and their
common area lies between these points: If the
curves y1 f1 ( x) and y2 f 2 ( x), where
f1 ( x) f 2 ( x) intersect in two points A( x a ) and
B ( x b) then, common area between the curves is
b b
( y1 y2 ) dx [ f1 ( x) f 2 ( x )] dx , as shown in
a a
Fig. 2.38 (a).
When two curves intersect at a point and the area
between them is bounded by x-axis: Area bounded
by the curves y1 f1 ( x), y2 f 2 ( x) and x axis is,
b
f1 ( x ) dx f 2 ( x ) dx , where, P( , ) is the
a
point of intersection of the two curves, as shown Fig.
2.38 (b). Figure 2.38: Area between the two curves
Positive and negative area: Area is always taken as
positive. If some part of the area lies above the x-axis and some part lies below x-axis, then
the area of two parts should be calculated separately and then add their numerical values to
get the desired area.
2
Example 2.225 [CE-1997 (1 mark)]: Area bounded by the curve y x and the lines x 4 and
y 0 is given by
(a) 64 (b) 64 3 (c) 128 3 (d) 128 4
2
Solution (b): Putting x 4 in y x we get y 16 ; so the line and
curve meet at (4,16) . So, the area of the shaded region bounded by
the curve y1 x 2 and the lines y2 0 and x 4 , given in figure, is
4 4 2
given as 0 ( y1 y2 )dx 0 ( x 0) dx ( x3 3) 04 4 3 3 64 3
Example 2.226 [ME-2003 (2 marks), ME-2012, PI-2012 (1 mark)]: The area enclosed between
2
parabola y x and the straight line y x is
(a) 1 8 (b) 1 6 (c) 1 3 (d) 1 2
2
Solution (b): Solving y x and y x we get x 0, y 0 and
x 1, y 1 ; so the line and curve meet at O (0, 0) and A(1,1) . Thus area of
the shaded region given in figure, is given as
1 1 2 1 1
0 ( y1 y2 )dx 0 ( x x ) dx ( x 2 2) ( x 3 3) (12 2) (13 3) 1 6
0 0
Example 2.227 [CE-2006 (2 marks)]: What is the area common to the circle r a and r 2a cos
?
(a) 0.524 a 2 (b) 0.614 a 2 (c) 1.047 a 2 (d) 1.228 a 2
Solution (d): The point of intersection of the two curves is a 2a cos cos 1 2
3, 5 3 . So the point of intersection of the two curves is (a, 3) and ( a, 5 3) . So the
5 3
area of common to the two curves r1 a and r2 2a cos is A (1 2) ( r22 r12 ) d
3
5 3 5 3
A (1 2)
3
(4a 2 cos2 a 2 ) d (1 2)
3
2a (1 cos 2 ) a d
2 2
5 3 5 3
A ( a 2 2) (1 2 cos 2 ) d ( a 2 2) ( sin 2 ) 3
1.228 a 2
3
2 2
Example 2.228 [ME-2009 (2 marks)]: Area enclosed between the curves y 4 x and x 4 y is
(a) 16 3 (b) 8 (c) 32 3 (d) 16
2 2
Solution (a): Solving y 4 x and x 4 y we get x 0, y 0 and
x 4, y 4 ; so the curves meet at O (0, 0) and A(4, 4) . Thus required area
4 4 4 4
is 0 ( y1 y2 )dx 0 2
x ( x 2 4) dx 2 x 3 2 (3 2) 0
(1 4)( x 3 3)
0
32 3
4(4 3) (1 4)(4 3) (32 3) (16 3) 16 3
[Similar question was also asked in MN-2014 (2 marks)]
Example 2.229: [CE-2016 (2 marks)]: The area of the region bounded by the parabola y x 2 1
and the straight line x y 3 is
(a) 59 6 (b) 9 2 (c) 10 3 (d) 7 6
Solution (b): The point of intersection of y x 2 1
and x y 3 is given as:
x 2 1 3 x x 2 x 2 0 ( x 2)( x 1) 0
x 2,1 . So the two curves meet at A(1, 2) and
B (2,5) . Thus required area of shaded region is given
1 2 1 2
by A {(3 x) ( x 1)}dx {2 x x }dx
2 2
1
x 2 x3 12 13 ( 2) 2 ( 2)3 7 20 27 9
A 2x 2 2( 2)
2 3 2 2 3 2 3 6 6 6 2
[Similar question was also asked in CE-2016 (2 marks)]
Example 2.230 [TF-2011 (1 mark)]: The area of ellipse with ‘ a ’ and ‘ b ’ as the length of major and
minor axis, respectively, is
(a) ab (b) ( a b ) 2 (c) ab 4 (d) ab 2
Solution (c): As the area of ellipse having ‘ A ’ and ‘ B ’ as length of semi-major and semi-minor axis
is AB ; so we have 2A a and 2B b , thus the required area of ellipse is AB ab 4 .
Surfaces of Solid of Revolution: The curved surface of the solid generated by the revolution,
about the x axis, of the area bounded by the curve y f ( x ) ,
the ordinates at x a , x b and the x axis is equal to
x b
2 y ds , as shown in Fig. 2.39.
xa
If the arc of the curve y f ( x ) revolves about y axis,
then the area of the surface of revolution (between proper
2
limits) 2 x ds, where ds 1 (dy dx) dx .
If the equation of the curve is given in the parametric form Figure 2.39: Area of surface of revolution
x f1 (t ) and y f 2 (t ) , and the curve revolves about x axis, then we get the area of the
t t2 t t2 t2
surface of revolution 2
t t1
yds 2
t t1
f 2 (t ) ds 2 f 2 (t )
t1
(dx dt ) 2
( dy dt ) 2 dt ,
Example 2.231 [AG-2008 (2 marks)]: Cycloid is formed by x a ( sin ) and y a(1 cos ) .
The surface area of the curved plane obtained from the rotation of the cycloid around x axis, from
0 to 2 , is
2 2 2 2
(a) 16 a 3 (b) 32 a 3 (c) 64 a 3 (d) 128 a 3
Solution (b): dx d a (1 cos ) and dy d a sin , so the surface area of the curved plane
from the rotation of the cycloid x a ( sin ) and y a(1 cos ) around x axis is
2 2
s 2
0
y (dx dt ) 2
( dy dt ) 2 d 2
0
a (1 cos ) {a (1 cos )}2 (a sin ) 2 d
2 2
s 2 a 2 2 sin 2 ( 2) 2(1 cos ) d 2 a 2 2 sin 2 ( 2) 4 sin 2 ( 2) d
0 0
2 2
s 2 a 2 4 sin 3 ( 2) d 2 a 2 {3sin( 2) sin(3 2)}d
0 0
2
s 2 a 2 6 cos( 2) (2 3) cos(3 2) 0 2 a 2 6 (2 3) 6 (2 3) (64 3) a 2
2 2
Example 2.232: The part of circle x y 9 in between y 0 and y 2 is revolved about y
axis. Find the volume of generating solid.
2 2
Solution: The part of circle x y 9 in between y 0 and y 2 is revolved about y axis.
Then a frustum of sphere will be formed. The volume of this frustum
2 2 2
x 2 dy (9 y 2 )dy 9 y (1 3) y 3 (46 3) cubic unit.
0 0 0
Volume of Solid of Revolution: If a plane curve is revolved about some axis in the plane of the
curve, then the body so generated is known as solid of revolution. The surface generated by the
perimeter of the curve is known as surface of revolution and the volume generated by the area is
called volume of revolution. For e.g. a right angled triangle when revolved about one of its sides
(forming the right angle) generates a right circular cones.
The volume of the solid generated by the revolution, about the x axis, of the area bounded by
b 2
the curve y f ( x) , the ordinates at x a , x b and the x axis is equal to a y dx , as shown
in Fig. 2.40.
The volume of the solid generated by the revolution, about y axis, of the area bounded by the
curve x f ( y ) , the abscissa y a and y b is given by (interchanging x and y in the above
b 2
formulae) a x dy .
If the equation of the generating curve be given by x f1 (t ) and y f 2 (t ) and it is revolved
b 2 t2
about x axis, then the formula corresponding to a y dx becomes t { f 2 ( t )}2 d { f1 ( t )} ,
1
The volume in the case when the above area is revolved about the line 2 is
(2 3) r 3 cos d .
If the generating
curve revolves
about any line AB
(which is different
from either of the
axes), then the
volume of
revolution is
2
( PN ) d (ON ) ,
Figure 2.41: Volume of the solid
as shown in Fig. Figure 2.40: Volume of the solid generated by generated by the revolution of the
2.41. the revolution of the curve about −axis curve about the line AB
Volume of an ellipsoid object having three semi axis whose length are a, b, c is V (4 3) abc .
Example 2.233 [ME-2010 (1 mark)]: The parabolic arc y x , 1 x 2 is revolved around the
x axis. The volume of the solid of revolution is
(a) 4 (b) 2 (c) 3 4 (d) 3 2
Solution (d): The required volume of the solid generated by the revolution about x axis, of the area
bounded by the curve y x , the ordinated at x 1 , x b and the x axis is equal to
2 2 2
V y 2 dx xdx ( x 2 2) ( 2)(2 2 1) 3 2
1 1 1
Example 2.234 [AG-2007 (2 marks)]: An ellipsoid object has three axes measuring 40 cm, 20 cm
and 20 cm respectively. The volume of the object is
(a) 4.23 litres (b) 8.38 litres (c) 12.63 litres (d) 17.05 litres
Solution (b): The length of semi axis is 20 cm, 10 cm and 10 cm. So
V (4 3) (20)(10)(10) 8.377 103 cm3 8.377 litres.
expression V R (1 h H ) dh
H 2 2
for
(a) 0 R (1 h H ) dr (b) 0 R (1 h H ) dh
0 H R 2
the volume of a cone is equal to
(c) (d) 0 2 rH (1 r R )dh 0 RH (1 r R ) dr
Solution (d): The given integral represents the volume of the cone, whose base is centred at the
origin, which is the integral of an infinite number of infinitesimally thin circular discs of thickness dh
and each of radius r , which is found by similar triangles r R ( H h ) H r R ( H h) H ; so
2
the surface area of the circular disc is A r 2 R 2 ( H h ) H .
2 2 2
H H h
2 H 2 h R H
So, volume V R dh 0 R 1 dh .
0
H H 3
Now, we consider options (a) and (d) only, because these contains
variable r , as the variable of integration. By integrating option (d), we
R R
get, V RH (1 r R ) 2 dr ( H R )( R 2 r 2 2 Rr )dr
0 0
2 R
V ( H R) R r (r 3) Rr 2 3
0
3
( H R ) R ( r 3) R
3 3
(1 3) R H .2
Example 2.236 [AE-2012 (2 marks)]: The volume of a solid generated by rotating the region
between semi-circle y 1 1 x 2 and straight line y 1 , about x axis, is
2 2 2 2
(a) (4 3) (b) 4 ( 3) (c) (3 4) (d) (3 4)
Solution (a): The curve y 1 1 x 2 is a circle with centre
(0,1) and radius 1, which is shown in figure. The required volume
of a solid generated by rotating the shaded region between semi-
circle y1 1 1 x 2 and straight line y2 1 , about x axis is the
volume generated by the region AOBCDA minus the volume
generated by DOCD, i.e.,
x 1 x 1 x 1
V y12 dx y22 dx ( y12 y22 ) dx
x 1 x 1 x 1
x 1 1
x 1 2 x 1
V (1 1 x 2 ) 2 dx ( 1 x 2 2 1 x 2 ) dx
x 1
x 1
V 2
x 1
x 0
( 1 x 2 2 1 x 2 ) dx 2 x ( x 3 3) 2 ( x 2) 1 x 2 (12 2) sin 1 ( x 1) x 0
3
V 2 1 (1 3) 2 0 (1 2)sin 1 0 2 (2 3) ( 2) ( 1
2
4) (4 3) .
Example 2.237 [XE-2016 (1 mark)]: The volume of the solid obtained by revolving the curve
y 2 x , 0 x 1 around y axis is
(a) (b) 2 (c) 2 (d) 5
Solution (d): The volume of the solid generated by the revolution, about y axis, of the area bounded
b 2
by the curve x f ( y ) , the abscissa y a and y b is given by V x dy . Here we have
a
b 1
0 x 1 0 y 1 ; and x f ( y ) y 2 . So V x 2 dy y 4 dy (1 5)( y 5 )10 5 .
a 0
Exercise: 2.5
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. (x
2
tan 1 x 3 ) (1 x 6 ) dx is equal to
1 3 1 3 2 1 3 2 1 3 2
(a) tan ( x ) c (b) (1 6)(tan x ) c (c) (1 2)(tan x ) c (d) (1 2)(tan x ) c
sin sin
1 1 1 1
2. x cos x x cos x dx is equal to
1 1
(a) (2 )[(2 x 1) sin x x (1 x )] x c (b) (2 )[(2 x 1) sin x x (1 x)] x c
1
(c) ( 2)[(2 x 1) sin x (1 x) ] x c (d) None of these
x
3. If e sin x dx (1 2)e x a c, then a
(a) sin x cos x (b) cos x sin x (c) tan x c (d) None of these
( x
2 4 2
4. 1) ( x x 1) dx (1 2) log k c , where k
x2 x 1 x2 x 1 x2 x 1 x2 x 1
(a) k (b) k (c) k (d) k
x2 x 1 x2 x 1 x2 x 1 x2 x 1
5. If 1 (1 sin x ) dx tan ( x 2) a b then
(a) a 4, b 3 (b) a 4 , b 3
(c) a 4 , b arbitrary constant (d) a 4 , b arbitrary constant
6. cos x (1 sin x) (2 sin x) dx
(a) log[(1 sin x ) (2 sin x)] c (b) log (2 sin x) (1 sin x) c
(c) log (1 sin x) (2 sin x) c (d) None of these
1 (sin cos ) d
3
7.
2 2
(a) log tan tan c (b) log tan (1 2) tan c
2
(c) log tan (1 2) tan c (d) None of these
2
3 1 x 1 x 1
8. 1 tan 2
x 1
tan
x
dx k , where k _____.
2
9. 0 sin x dx k , where k _____.
e2
10. e (loge x) x dx _____.
1
2
11. 0 1 (1 sin x) dx _____.
1
12. Suppose f is such that f ( x ) f ( x) for every real x and 0 f ( x) dx 5 , then
0
1 f (t ) dt k , where k _____.
2 2
13. 2 1 x dx _____.
1.5 2
14. 0 [ x ]dx , where [] denotes the greatest integer function.
(a) 2 2 (b) 2 2 (c) 1 2 (d) 2 1
cos x
e sin x , x 2 3
15. If f ( x ) , then 2 f ( x)dx _____.
2, otherwise
sin 2 x 3
16. 0 e cos xdx _____.
12
17. 1 2 [ x] ln (1 x) (1 x) dx k , where [] denotes greatest integer function, then k ____.
1
18. 1 log x x 2 1 dx _____.
2
19. (1 x ) sin x cos 2 xdx _____.
cos 2 x
20. If n is any integer, then 0 e cos3 (2n 1) x dx _____.
3 4
21. 4 1 (1 cos x) dx _____.
22. If 0 x f (sin x ) dx k f (sin x) dx , then the value of k _____.
0
2 x sin 2 n x
23. For n 0, dx
0 sin 2 n x cos 2 n x
(a) 2 (b) 2 2 (c) 3 2 (d) 4 2
a T
24. If f ( x ) is a continuous periodic function with period T , then the integral I f ( x ) dx
a
(a) 2a (b) 3a (c) Independent of a (d) None of these
1n
25. lim{(n !) n} (a) e (b) e 1
(c) 1 (d) None of these
n
2 4 6
26. 0 sin x cos xdx 3 k , where k _____.
2 7
27. 0 sin xdx _____.
x 2 2
28. Let f : (0, ) R and f ( x ) f (t ) dt . If f ( x ) x (1 x) then f (4) _____.
0
x 2
29. Let f ( x ) 2 t 2 dt . Then the real roots of the equation x f ( x) 0 are
1
In addition, if the vector function is defined at t0 and lim F(t ) F(t0 ) then F (t ) is said to be
t t0
continuous at t0 . A vector function F (t ) that is continuous for each t in the interval a t b is said
to be continuous over the interval. A vector function of a single real variable that is not continuous at
a point t0 is said to be discontinuous at t0 .
Let u(t ) and v (t ) be differentiable functions of t over some interval a t b , with C an arbitrary
constant vector and c be an arbitrary constant scalar. Then rules for differentiation of vector functions
of a single real variable over the interval a t b are:
( f u ) f u f u , where f f ( ) (u v ) u v u v
( u v ) u v , where , are constants (u v) u v u v
(u v w ) u v w u v w u v w d du df ( )
u f ( )
[u ( v w )] u ( v w ) u ( v w ) u ( v w ) d df d
Differentiation of a constant vector: C 0
If u(t ) is a differentiable function of t and t t ( s) is a differentiable function of s , then
du du dt du du1 dt du2 dt du3 dt
; or if u (t ) u1 (t )i u2 (t ) j u3 (t )k then, i j k
ds dt ds ds dt ds dt ds dt ds
Example 2.238 [MN-2012 (2 Solution (a): R t 2 iˆ 2t ˆj d R dt 2t iˆ 2t ˆj . Thus slope
marks)]: The angle between the
2
ˆ ˆ
tangents to the curve R t i 2t j of curve at t 1 is m 1 d R dt 2 iˆ 2 ˆj and
t 1
at the point t 1 is m2 d R dt 2 iˆ 2 ˆj . So, the angle between the
t 1
(a) 2 (b) 3
tangents at the points t 1 is
(c) 4 (d) 6
cos m1 m2 m1 m2 ( 4 4) m1 m2 0 2
Vector Differential: The vector differential dF for the function given in Eq. 2.1 is defined as,
dF ( df1 dt )i (df 2 dt ) j (df 3 dt )k dt (2.3)
If s is the arc length measured from some fixed point, along space curve , which is defined by the
position vector r (t ) x1 (t ) i x2 (t ) j x3 (t ) k for t t0 , then dx1 ( dx1 dt ) dt , dx2 ( dx2 dt ) dt and
dx3 ( dx3 dt ) dt . As shown in Fig. 2.43, the differential element of arc length ds along is
12 12
dx1 2 dx2 2 dx3 2 2 2 2
ds dr dx1 dx2 dx3
ds dt (2.4)
dt dt dt dt dt dt dt dt
When t is the time and r (t ) is a position vector in space, ds dt dr dt is the speed with which
the tip of position vector r (t ) traces out a space curve .
The unit tangent vector T along as a function of t is given by T ( dr dt ) dr dt . Also,
ds dt dr dt dr dt ( ds dt )T .
Since T is the unit tangent vector along T T 1 , so
differentiating it w.r.t. s (dT ds ) T T (dT ds ) 0
T ( dT ds ) 0 (since dot product is commutative) which
shows that T and d T ds are orthonormal. The unit vector
N in the direction of d T ds at a point r r ( s) on is
called the principal normal to at r ( s ) , and so
N (dT ds ) ( dT ds ) for dT ds 0 .
When the relation between dT ds and N at a point
r r ( s) on is written in the form dT ds k ( s ) N , the
Figure 2.43: The geometrical relationship
k ( s ) is called the curvature of the curve at r r ( s) , and between the differentials , , ,
( s ) 1 k ( s ) is called the radius of curvature of the curve and
at r r ( s) . As N is a unit vector, so k ( s) dT ds . In case of a smooth plane curve , the
circle of curvature at a point P on is tangent to at P with radius ( s ) 1 k ( s ) , and such
that its centre lies on the concave side of .
If the curvature is required in terms of the parameter t , the relationship between k ( s ) and k (t ) is
dT dT ds dT ds dT dr
given from the chain rule as: k (t ) k (t )
dt ds dt dt dt dt dt
The unit binormal vector to the curve at r r ( s) is defined as, B T N
Three unit vectors T , N and B at a point r r ( s) on the space curve form a triad of
mutually orthogonal unit vectors whose orientation depends on the location of the point on the
curve .
grad T (1,1,1)
2 i 2 j 4 k 22 22 42 24 2 6
[Similar questions were also asked in ME-1998, EC-2014 (1 mark), EE-2005, MT-2014 (2
marks)]
direction is T y ( y ) 400 50 ze (t x
2
y2 )
50 ze (t x2 y2 )
( 2 y ) 100 yze (t x
2
y2 )
Directional Derivatives and its properties: The rate of change of a function f ( x, y, z ) in the
direction of the unit vector vˆ l i m j n k , is called the directional derivative and is denoted by
Dvˆ f . Thus Dvˆ f ( x, y , z ) lim f ( x lh, y mh, z nh ) f ( x, y , z ) h . The directional derivative
h 0
Dvˆ f can be interpreted in terms of the ordinary operation of differentiation by considering Fig. 2.44.
In the diagram, a straight line T in space in the direction of a given
vector v̂ passes through a fixed point P and Q is a general point on
the line T at a distance s from P . The direction derivative Dv f is
then given by Dvˆ f df dv lim f (Q ) f ( P ) s . Let v be a unit
s 0
vector defined as vˆ l i m j n k , where l , m and n are the Figure 2.44: The directional
direction cosines of the tangent to the space curve , so we have derivative
2 2 2 12
dx dr dy dr dz dr dr dx dy dz
l , m and n with . Then as the
dt dt dt dt dt dt dt dt dt dt
scalar product of a vector F and the unit vector v̂ is the projection of F in the direction v̂ , it
follows that, directional derivative of f in the direction of v̂ , is given as,
Dvˆ f vˆ grad f l (f x ) m(f y ) n(f z ) (2.6)
The most rapid increase of a differentiable function f ( x, y, z ) at a point P in space occurs in the
direction of the vector vˆ P grad f ( P ) . The directional derivative at P is then given by
2 12
2 2
Dvˆ f ( P ) grad f ( P ) (f x ) P (f y ) P ( z ) P .
The most rapid decrease of a differentiable function f ( x, y, z ) at a point P in space occurs when
the vector vˆ P just defined in previous point and grad f are oppositely directed, so that
vˆ P grad f ( P ) . The directional derivative at P is then the negative of the result in the
2 12
2
previous point and hence, Dvˆ f ( P ) grad f ( P ) (f x) P (f y ) P ( z ) P
2
.
There is a zero local rate of change of a differentiable function f ( x, y, z ) at a point P in space in
the direction of any vector vˆ P that is orthogonal to grad f at P , so that vˆ P grad f ( P ) 0 .
When a scalar function f defined over a region D of space is suitably differentiable, the vector-
valued function grad f defines a vector field over D in terms of the scalar field defined by f .
Example 2.243 [CE-2009 (2 mark)]: For a scalar function f ( x, y , z ) x 2 3 y 2 2 z 2 , the
directional derivative at the point P (1, 2, 1) in the direction of a vector i j 2k is
(a) 18 (b) 3 6 (c) 3 6 (d) 18
2 2 2
Solution (c): As the scalar point function is f ( x, y , z ) x 3 y 2 z
grad f f (f x) i (f y ) j (f z ) k 2 x i 6 y j 4 z k grad f (1,2, 1)
2 i 12 j 4 k
. So the directional derivatives of f at (1, 2, 1) in the direction of v i j 2 k is
v (i j 2 k ) 2 12 8 18
Dvˆ f vˆ grad f grad f (2 i 12 j 4 k ) 3 6
v i j 2k 12 ( 1) 2 22 6
[Similar questions were also asked in CE-2002, CE-2006, TF-2009, ME-2008 (2 marks), XE-
2008 (1 mark)]
Example 2.244 [PI-2007 (1 mark)]: f ( x) x is a function defined for real numbers x . The
directional derivative of f at x 0 in the direction d 1 is
(a) 1 (b) 0 (c) 1 2 (d) –1
Solution (a): As the given function is single variable function so, direction derivative of f ( x) x at
f (0 1 h ) f (0) 0h 0
x 0 in the direction d 1 is Dd f ( x ) x 0 lim lim lim1 1 .
h0 h h 0 h h 0
Example 2.245 [CH-2011 (2 marks)]: Unit vectors in x and z directions are i and k respectively.
Which ONE of the following is the directional derivative of the function F ( x, z ) ln( x 2 z 2 ) at the
point P : (4, 0) , in the direction of (i k ) ?
(a) (1 2 2 )i (b) i (c) 1 (d) 1 2 2
2 2
Solution (d): As the scalar point function is F ( x, z ) ln( x z )
grad F F (F x) i (F y ) j (F z ) k 2 x ( x 2 z 2 ) i 2 z ( x 2 z 2 ) k
grad F (4,0 )
(1 2) i . So the directional derivatives of F at (4, 0) in the direction of v i k is
(i k ) 1
v 12 1
Dvˆ f vˆ grad f grad f
i
v ik 2 12 12 2 2
[Similar question was also asked in CH-2007 (2 marks)]
Example 2.248 [ME-2000 (2 marks)]: The maximum value of the directional derivative of the
function 2 x 2 3 y 2 5 z 2 at a point (1,1, 1) is
(a) 10 (b) –4 (c) 152 (d) 152
Solution (c): As grad P
P
( x ) iˆ ( y ) ˆj ( z ) kˆ P
is the direction of
Example 2.249 [XE-2011 (1 mark)]: Consider the function f ( x, y , z ) x 3 e y sin z and the point
P 1, 0, 2 . The value of f DOES NOT change due to small displacement of P along the direction
of
(a) 1, 0, 2 (b) 1, 1,1 (c) 1, 3, 0 (d) 2, 0, 1
Solution (c): Let we have a unit vector vˆ l i m j n k , then directional derivative of
f ( x, y , z ) x 3 e y sin z at point P 1, 0, 2 along vˆ l i m j n k must be zero when the value of
f does not change due to small displacement of P along the direction of v̂ ; i.e.,
Dvˆ f P
0 f P
vˆ 0 (f x ) i (f y ) j (f z )k P (l i m j n k ) 0
( x )( x 3e y sin z )i ( y )( x 3e y sin z ) j ( z )( x 3 e y sin z )k P
(l i m j n k ) 0
(3 x e 2 y
sin z )i ( x 3e y sin z ) j ( x3 e y cos z )k (l i m j n k ) 0
1,0, 2
Rules for the gradient operator: Let the gradients of f and g be defined over a region D , then
grad (cf ) c grad f c f , where c is a scalar constant.
grad ( f g ) grad f grad g f g
grad ( f g ) ( f g ) f grad g g grad f f g g f
grad f g g grad f f grad g g 2 g f f g g 2 , g 0
grad f ( g ) ( df dg )g
The above results are obtained by the usual rules for partial differentiation to each component of the
gradient function on the left, and then recombining the results to obtain the expression on the right.
Let us prove the third result, grad ( f g ) ( fg ) x i ( fg ) y j ( fg ) z k
grad ( f g ) f (g x ) g (f x ) i f (g y ) g (f y ) j f (g z ) g (f z ) k
grad ( f g ) f (g x)i (g y ) j (g z )k g (f x )i (f y ) j (f z )k
grad ( f g ) f grad g g grad f f g g f
Gradient as a Surface Normal Vector and Tangent plane to the surface : Let S be a
surface represented by f ( x, y, z ) c , where c is any
constant and f is differentiable function. Such a surface
is called a level surface of f , and for different c we get
different level surfaces. Now let C be a curve on S
through a point P of S . As a curve in space, C has a
representation r (t ) x (t ) i y (t ) j z (t ) k . For C to lie
on the surface S , the components of r (t ) must satisfy
f ( x, y, z ) c 0 . Now a tangent vector of C is
dr (t ) dt r (t ) x(t ), y (t ), z(t ) and the tangent
vector of all curves on S passing through P will Figure 2.45: Gradient as a Surface Normal Vector
generally form a plane, called the tangent plane of S at P (exceptions occur at the edges of S ). The
normal of this plane (the straight line through P perpendicular to the tangent plane) is called the
surface normal to S at P . A vector in the direction of the surface normal is called a surface normal
vector of S at P , which is obtained by differentiating f ( x, y, z ) c 0 w.r.t. t with the help of
chain rule as, (f x)( dx dt ) (f y)( dy dt ) (f z )(dz dt ) (grad f ) r 0 . Hence grad f is
orthogonal to all the vectors r in the tangent plane, so that it is a normal vector of S at P .
The vector equation of a plane with normal n containing the point P0 with position vector r0 is
(r r0 ) n 0 , where r is the position vector of an arbitrary point on the plane. If we set
r x i y j z k and r0 x0 i y0 j z0 k and identity n with grad f at P0 , where
grad f ( P0 ) (f x) P0 i (f y ) P0 j (f z ) P0 k .
The required tangent plane to the surface at P0 ( x0 , y0 , z0 ) is to be given by
( x x0 )(f x ) P0 ( y y0 )(f y ) P0 ( z z0 )(f z ) P0 0 .
Example 2.250 [XE-2007 (2 marks)]: The maximum value of the (a) 2 (b) 3
function 2 x 3 y 4 z on the ellipsoid 2 x 2 3 y 2 4 z 2 1 is (c) 6 (d) 9
Solution (b): Let f ( x, y , z ) 2 x 3 y 4 z and g ( x, y, z ) 2 x 2 3 y 2 4 z 2 1 , so we have to
maximize f ( x, y, z ) subject to g ( x, y , z ) . Using Lagrange multiplier , we have f g
2 i 3 j 4 k (4 x i 6 y j 8 z k ) x y z 1 2 . Putting values of x, y, z in
g ( x, y , z ) 0 , 2 (4 ) 3 (4 ) 4 (4 ) 1 0 9 (4 ) 1 3 2 . So critical
2 2 2 2
A more general way of defining the divergence of a vector field F (r) at the point r is to
enclose the point in an elementary volume V and find the flow or flux out of V per unit
flow out of V
volume. Thus, div F F lim .
V 0V
A non-zero divergence at a point in a fluid measures the rate, per unit volume, at which the fluid is
flowing away from or towards that point. That implies that either the density of the fluid is changing
at the point or there is a source or sink of fluid there. When div F 0 , there is a net flow of F out of
the volume; and when div F 0 , there is a net flow of F into the volume. In the case of a non-
material vector field, for example temperature gradient in heat transfer, a nonzero divergence
indicates a point of generation or absorption. When the divergence is everywhere zero, the flow
entering any element of the space is exactly balanced by the outflow. This implies that the lines of
flow of the field F (r ) where div F 0 must either form closed curves or finish at boundaries or
extend to infinity. Vectors satisfying this condition are sometimes termed solenoidal.
Properties of the divergence operator: Let the vector fields F , G and the scalar fields , be a
suitably differentiable; let a , b are constants. Then, divergence operator has the following properties:
div ( aF ) a div F a F div (aF bG ) a div F b div G a F b G
div (grad ) 2 div ( ) 2 grad grad 2
div ( F) div F F F F
div ( ) div ( ) 2 2
Divergence of a vector in Cylindrical and Spherical Coordinates
Cylindrical Polar 1 (r u r ) 1 u u z
: u , where x r cos , y r sin and z z
Coordinates: r r r z
1 ( r 2u r ) 1 (u sin ) 1 u
Spherical u 2
, where
: r r r sin r sin
Coordinates:
x r sin cos , y r sin sin , z r cos
Example 2.251 [ME-2001, EE-2010 (1 mark), MT-2014 (2 marks)]: Divergence of the three-
dimensional radial vector field r is
(a) 3 (b) 1 r (c) iˆ ˆj kˆ (d) 3(iˆ ˆj kˆ)
Solution (a): As any 3-D radial vector field r can be represented as r x iˆ y ˆj z kˆ . So
div r r ( x ) i ( y ) j ( z )k ( x iˆ y ˆj z kˆ) (x x ) (y y ) (z z ) 3
Example 2.252 [CE-2007 (2 marks)]: A velocity vector is given as (a) 9 (b) 10
2
2
V 5 xyi 2 y j 3 yz k . The divergence of this velocity vector at (1,1,1) is (c) 14 (d) 15
2
2
Solution (d): div V V ( x) i ( y ) j ( z )k (5 xyi 2 y j 3 yz k )
div V ( x )(5 xy ) ( y )(2 y 2 ) ( z )(3 yz 2 ) 5 y 4 y 6 yz div V 5 4 6 15
(1,1,1)
[Similar questions were also asked in EC-1998 (5 marks), MN-2009, MT-2007, TF-2010 (2
marks) (2 marks), ME-2008, ME-2009, EC-2013 (1 mark)]
Example 2.253 [MN-2014 (1 mark)]: The divergence of the vector b ( x y )( y iˆ x ˆj ) is
(a) y x (b) x y (c) x 2 y 2 (d) y 2 x 2
Solution (b): div b b ( x ) i ( y ) j ( z )k ( yx y 2 ) i ( x 2 xy ) j 0 k
2 2
div b ( x )( yx y ) ( y )( x xy ) ( z )(0) y x
[Similar question was also asked in ME-2014 (1 mark)]
Example 2.254 [IN-2012, EE-2012, EC-2012 (2 marks)]: The direction of vector A is radially
n
outward from the origin, with A kr where r 2 x 2 y 2 z 2 and k is constant. The value of n
for which A 0 is
(a) –2 (b) 2 (c) 1 (d) 0
2 2 2 2
Solution (a): As r x y z is the equation of a sphere of radius r and thus in spherical
1 2 1 2 n k n 2
coordinates, for u u 0 , A 2 ( r ur ) 2 (r kr ) 2 ( r ) k ( n 2)r n 1 . Thus
r r r r r r
for A 0 , we must have n 2 0 n 2 .
(a) 2 z cos z 2 (b) sin xy 2 z cos z 2 (c) x sin xy cos z (d) None of these
Solution (a):
div V V {( x )iˆ ( y ) ˆj ( z ) kˆ} {( x cos xy y )iˆ ( y cos xy ) ˆj (sin z x y ) kˆ}
2 2 2
div V V ( x )( x cos xy y )iˆ ( y )( y cos xy ) ˆj ( z )(sin z x y ) kˆ
2 2 2
2 2
div V V cos xy xy sin xy cos xy xy sin xy 2 z cos z 2 z cos z
[Similar question was also asked in AG-2011 (2 marks)]
Example 2.256 [EC-2014 (2 marks)]: If r xaˆ x yaˆ y zaˆ z and r r , then div{r 2(ln r )} …
Solution: In spherical coordinates for ˆ ˆ 0 , f (f r ) rˆ ; so for f ln r ,
2 2
(ln r ) ( r )(ln r ) 1 r div{r (ln r )} div( r r ) div( r ) ( x ) x ( y ) y ( z ) z 3
Example 2.257 [EE-2014 (1 mark)]: Let ( f v ) x 2 y y 2 z z 2 x , where f and v are scalar and
vector fields respectively. If v yiˆ zjˆ xkˆ , then v f is
(a) x 2 y y 2 z z 2 x (b) 2 xy 2 yz 2 zx (c) x y z (d) 0
Solution (a): ( f v ) div ( f v ) f div v v f f ( v ) v f .
2 2 2
v ( x) y ( y ) z ( z ) x 0 v f ( f v ) x y y z z x
ˆ ˆ ˆ
j k ( x y )iˆ ( z 2 xy ) ˆj ( xy ) kˆ
2
Solution: div V V i
x y z
div V V ( x2 y) ( z 2 xy ) ( xy ) 2 x 2 x 0 0
x y z
Example 2.259 [MT-2016 (1 marks)]: If V x 2 y iˆ y 2 x ˆj xyz kˆ , the divergence of V is
(a) x 3 y y 3 x xyz 2 (b) x 2 y y 2 x xyz (c) 5xy (d) 0
ˆ ˆ
Solution (c): Divergence of V is V , where iˆ j k.
x y z
ˆ ˆ 2 ˆ 2 2
j k ( x y i y x ˆj xyz kˆ ) ( x y ) ( y x ) ( xyz )
2
So V iˆ
x y z x y z
V 2 xy 2 xy xy 5 xy .
Curl of a Vector: The operation, F , is called the curl of a vector and it associates a vector
function with a differentiable vector field F . If F F1 i F2 j F3 k is a differentiable vector field,
the curl of F , written as curl F , is the scalar function defined in terms of Cartesian coordinates as,
i j k
F F F F F F
curl F F x y z 3 2 i 1 3 j 2 1 k (2.9)
y z z x x y
F1 F2 F3
[Equation 2.9 was asked in AG-2016 (2 marks)]. More generally, the component of the curl of a
vector field F (r ) in the direction of the unit vector n̂ at a point L is found by enclosing L by an
elementary area S that is perpendicular to L , as shown in Fig. 2.46, and
calculating the flow around S per unit area. Thus,
(curl F) nˆ lim (flow around S ) S .
S 0
Properties of curl and other operators: Let F and G be vector functions and
let be a scalar function, all of which are suitable differentiable. The operators
grad, div, and curl can be combined in various ways that lead to identities, the Figure 2.46:
results which are listed as follows: Circulation around
curl (grad ) 0 , where 0 is vector of zero magnitude [This point the element ∆
was asked in ME-1996, EE-2013, MT-2014 (1 mark)]
div (curl F) ( ) 0 grad ( F ) ( F) F F
curl ( F ) (curl F ) F (grad ) ( F) F ( ) [This point was asked in ME-1995 (1
mark)]
grad (F G ) F (curl G ) G (curl F) (F )G (G )F
div (F G ) (F G ) G (curl F ) F (curl G ) G ( F ) F ( G )
curl (F G ) (F G ) F (div G ) G (div F ) (G )F (F )G
F ( G ) G ( F ) (G )F (F )G
curl (curl F) ( F ) ( F) 2 F grad (div F ) 2 F [This point was asked in EC-
2006 (1 mark)]
grad (div F ) ( F ) 2 F
Example 2.261 [ME-2003 (2 marks)]: The vector field F x i y j (where i and j are unit
vectors) is:
(a) divergence free, but not irrotational (b) irrotational, but not divergence free
(c) divergence free and irrotational (d) neither divergence free nor irrotational
Solution (c): div F F i j k x i y j 0 k x ( y ) 1 1 0 .
x y z x y
i j k
So the vector field F is divergence free. As, curl F F x y z 0 i 0 j 0 k . So
x y 0
the vector field F is irrotational.
The table above lists formulas for the vector derivatives of some common expressions. The derivative
f xT is the transpose of f x . In the table x is an n vector, a is a constant scalar, b is a
constant vector; and A is a constant matrix.
Example 2.267 [IN-2014 (2 marks)]: A scalar valued function is defined as f ( x) xT Ax bT x c ,
where A is a symmetric positive definite matrix with dimension n n ; b and x are vectors of
dimension n 1 . The minimum value of f ( x ) will occur when x equals
(a) ( AT A) 1 b (b) ( AT A) 1 b (c) ( A1b ) 2 (d) ( A1b ) 2
Solution (c): For extremum value of the given function we have to differentiate the given function
T
w.r.t. the vector x , f (x ) ( x Ax) (bT x) (c) ( A AT ) x b 0 2 Ax b (as A
x x x x
T
is symmetric matrix so A A ) and now equating it to 0 for getting the critical point, i.e.,
( x) f ( x) 0 2 Ax b 0 2 A1 Ax A1b x (1 2) A1b . Now we have to check,
whether x (1 2) A1b is the point of maxima or minima; so
2 2
f ( x) (2 Ax) (b) 2 A 0 2 A f ( x) 0 (as A is a positive matrix).
x 2 x x x 2 x (1 2) A1b
None
(a) (b) (c) (d) of
these
3. The region in which the pressure field f ( x, y ) 9 x 2 4 y 2 varies between 36 and 144 is between
(a) two ellipses (b) two parabolas (c) two hyperbolas (d) two straight lines
4. The parametric representation of the straight line through a point A(3,1,5) in the direction of a
vector b 4 i 7 j k is
(a) (5 t ) i (3 4t ) j (1 7t ) k (b) (3 4t ) i (5 t ) j (1 7t ) k
(c) (3 4t ) i (1 7t ) j (5 t ) k (d) 3 i j 5 k
5. The curve given by the parametric representation r (t ) t i (t 3 2) j 0 k is
(a) y x 2 2 , z 0 (b) y x 3 2 , z 0 (c) y x 2 2 , z 2 (d) y x 3 6 , z 2
6. For a curve given by r (t ) t i t 3 j , the unit tangent vector u(t ) at a point P (1,1, 0) is
(a) i 3 j (b) (1 10) i (3 10) j (c) i 3 k (d) (1 10) i (3 10) k
32
7. The length of the curve given by r (t ) t i t j from (0, 0, 0) to (4,8, 0) is _____.
8. The magnitude of tangential acceleration for the curve given by r (t ) 5t 2 k is _____.
9. If w ( x 2 y 2 z 2 ) 1 2 , x cos t , y sin t , z t then dw dt at t 1 is _____.
10. Which one of the following depends on the choice of a particular coordinate system?
(a) Gradient (b) Divergence (c) Gradient and Curl (d) None of these
11. The directional derivative of a scalar point function is a function of
(a) position (b) direction (c) both (a) and (b) (d) either (a) or (b)
2 2
12. The grad f for the function f ln( x y ) at (2, 0) is
(a) i (b) j (c) i j (d) i j
13. The direction of maximum decrease of flow of heat in a temperature field T x y at P (8, 1) is
(a) 8 i j (b) 8 i j (c) i 8 j (d) i 8 j
14. The unit normal vector for the surface z ( x 2 y 2 ) 1 2 at the point P (6,8,10) is parallel to
(a) 6 i 8 j 10 k (b) 8 i 6 j 10 k (c) 8 i 6 j 10 k (d) 6 i 8 j 10 k
15. For the vector field v 2 x i 4 y j 8 z k , the potential field f such that v grad f is
(a) 2 x 2 y 2 4 z 2 (b) x 2 y 2 4 z 2 (c) x 2 2 y 2 4 z 2 (d) x 2 y 2 z 2
16. For the vector field v xy i 2 xy j , the potential field f such that v grad f is
(a) xy (b) x 2 y 2 (c) 2x 2 y 2 (d) f has no potential
17. The magnitude of directional derivative of f ( x 2 y 2 z 2 ) 1 2 at the point P(3, 0, 4) in the
direction of a i j k is _____.
18. The magnitude of directional derivative of f e x cos y at the point P (2, , 0) in the direction of
a 2 i 3 j is _____.
19. The divergence of vector field given by v e x (cos y i sin y j) is
(a) 2e x cos y (b) 2e x sin y (c) e x cos y (d) e x sin y
20. The divergence of vector field given by v ( x 2 y 2 ) 1 ( y i x j) is _____.
21. Consider the flow with velocity vector v y i . This flow is
(a) Compressible and Irrotational (b) Incompressible and Irrotational
(c) Incompressible and not Irrotational (d) Compressible and not Irrrotational
22. Which of the following is correct?
(a) div ( ) div ( ) (b) curl (curl F ) grad (div F ) F
(c) div ( ) div ( ) 2 2 (d) curl (curl F) div (grad F ) 2 F
23. If f tan 1 ( y x ) , then 2 f is _____.
24. If the vector field is given by v xyz ( x i y j z k ) , then curl v at point P (1,1,1) is
(a) i j k (b) i j (c) j k (d) 0
25. If u y i z j x k and v yz i zx j xy k , then u curl v
(a) i j k (b) i j (c) j k (d) 0
26. The necessary and sufficient condition that the force field F is conservative is
(a) grad F 1 (b) curl F 0 (c) div F 0 (d) curl F 1
27. The magnitude of maximum directional derivative of f x 2 y xyz z 3 x in a direction from the
point P (1, 0,1) is _____.
28. The maximum value of the directional derivative of f x 3 y 2 z 2 from P (1, 1,1) is _____.
29. If velocity vector of a flow is v ( x 2 y az ) i (bx 3 y 2 z ) j (3 x cy 2 z ) k , then the
value of a, b, c such that the flow is irrotational is, respectively, _____.
(a) 2, 3, 2 (b) 3, 2, 2 (c) 2, 2, 3 (d) None of these
30. The angle ( ) between the two surfaces x 2 y 2 z 16 and x 2 2 y z 2 9 at P (1,1, 2) is
(a) 0 90o (b) 90o 180o (c) 0o (d) 90o
31. The relationship between a, b, c such that f ax 2 by 2 cz 2 satisfies Laplace equation is
(a) a b c 0 (b) a b c 0 (c) a b c 0 (d) a b c 0
32. Which one of the following satisfies Laplace’s equation?
(a) f 2 x 2 y 2 z 2 (b) f x 2 2 y 2 3 z 2 (c) f x 2 3 y 2 2 z 2 (d) None of these
Example 2.268 [EC-2008 (2 marks)]: The value of the integral of the function g ( x, y ) 4 x 3 10 y 4
along the straight line segment from the point (0, 0) to the point (1, 2) in the x y plane is
(a) 33 (b) 35 (c) 40 (d) 56
Solution (a): The equation of straight line segment from the point (0, 0) to the point (1, 2) in the
3 4 3 4
x y plane is y 0 {(2 0) (1 0)}( x 0) y 2 x g ( x, y ) 4 x 10(2 x) 4 x 160 x .
x 1 x 1
Thus x 0 g ( x, y )dx x 0 (4 x
3
160 x 4 ) dx {x 4 160(1 5) x 5 }xx 10 1 (160 5) 0 33 .
1
, 0 t 1 , the value of the integral 0 F {r (t )} (d r dt ) dt is
(a) 7 12 (b) 17 12 (c) 27 12 (d) 37 12
Solution (d): F {r (t )} 5t iˆ t ˆj t kˆ and d r dt iˆ ˆj kˆ thus
2 3
1 1 1
0 F {r (t )} (d r dt ) dt (5t iˆ t 2 ˆj t 3 kˆ ) (iˆ ˆj kˆ ) dt (5t t 2 t 3 )dt
0 0
4 t 1
2
I (5 2)t (1 3)t (1 4)t
3
t 0
37 12
Example 2.271 [EE-2009 (2 marks)]: F ( x, y ) ( x 2 xy )aˆ x ( y 2 xy )aˆ y . Its line integral over the
straight line from ( x, y ) (0, 2) to ( x, y ) (2, 0) evaluates to
(a) –8 (b) 4 (c) 8 (d) 0
Solution (d): The curve is a straight line whose equation is y 2 (0 2) (2 0) ( x 0)
y x 2 , on which any point can be taken as r (t ) t aˆ x ( t 2) aˆ y r (t ) aˆ x aˆ y with
0 t 2 . So, F r (t ) r (t ) t 2
t ( t 2) aˆ x (t 2) 2 t ( t 2) aˆ y {aˆ x aˆ y }
F r (t ) r(t ) t 2 t (t 2) (t 2) 2 t (t 2) 2t 2 4t . Thus, line integral of F ( x, y ) over a
2 2 2 2
straight line is C F (r ) dr 0 F r2 (t ) r2 (t )dt 0 (4t 4)dt (2t 4t ) 8 8 0
0
[Similar question was also asked in PI-2009 (2 marks)]
Example 2.272
[EC-2010 (2 (a) 0
2
marks)]: If A xyaˆ x x aˆ y then (b) 2 3
(c) 1
C A dl over the path shown in
(d) 2 3
the figure is
Solution (c): We have xy plane so dl dx aˆ x dy aˆ y A dl xy dx x 2 dy . So,
C A d l y 1 A d l x 2 3
Adl
y 3
Adl
x 1 3
Adl . Along y 1, A dl x dx ;
x2 3 , A dl (4 3) dy ; y 3 , A dl 3 x dx ; x 1 3 , A d l (1 3) dy . Thus,
x2 3 y 3 4 x 1 3 y 1 1 x
2 x2 3
4
y 3
3x
2 x 1 3
1
y 1
C A d l x 1 3
x dx
y 1 3
dy
x 2 3
3 x dx
y 3 3
dy
2
3
y
y 1
2
3
y
y 3
x 1 3 x2 3
C A d l (1 2) (4 3) (1 3) (4 3)(3 1) (3 2) (1 3) (4 3) (1 3)(1 3) 1
2 2
Example 2.273 [EE-2013 (1 mark)]: Given a vector field F y xax yza y x az , the line integral
x 0 x 0
F dr
AB x 1
( x 3 3x 2 4 x 1)dx (1 4) x 4 x 3 2 x x 1 4 . Similarly, along
2
x 1
BO ,
y 0 y 0
x 0 dx 0 F d r y 3 dy y 3dy (1 4) y 4 1 4 . So, the work done by the
BO BO y 1 y 1
force F ( x x 2 )iˆ ( x 2 y 3 ) ˆj in moving a particle along the given path is
OABO F d r OA F d r AB F d r BO F d r (5 6) (1 4) (1 4) 1 3
[Similar question was also asked in CH-2007 (2 marks)]
Example 2.277 [XE-2014 (2 marks)]: If the work done in moving a particle once around a circle
x 2 y 2 4 under the force field F ( x, y ) (2 x ay ) iˆ (2 y ax) ˆj is 16 , then a is equal to ……
Solution: Any point on x 2 y 2 2 2 can be taken as x 2cos , y 2 sin with 0 2 ,
r 2 cos i 2 sin j k ( dr dt ) r (t ) 2 sin i 2 cos j . So the line integral
2
I F r ( ) r ( )d (4 cos 2a sin ) i (4 sin 2a cos ) j ( 2 sin i 2 cos j)d
C 0
2 2 2
I
0
(8sin cos 4a sin 2
) (8sin cos 4 a cos 2 ) d 4 ad 4 a
0 0
8a
I 8a 16 a 2 .
Example 2.278 [MA-2016 (2 marks)]: Let be the triangular path connecting the points (0, 0) ,
(2, 2) and (0, 2) in the counter-clockwise direction in 2 . Then I 3
{sin( x )dx 6 xydy} is equal
to _____.
Solution: In counter clock wise direction the given line integral I 3
{sin( x )dx 6 xydy} can be
3 3 3
written as I {sin( x ) dx 6 xydy} {sin( x ) dx 6 xydy} {sin( x ) dx 6 xydy} . Now
OB BA AO
3 x 2
along OB : y x dy dx , so I1 {sin( x ) dx 6 xydy} {sin( x3 ) 6 x 2 }dx
OB x 0
x2 x 2 x 2 x 2
I1 sin( x3 ) dx 6 x 2 dx sin( x3 ) dx 2( x3 ) 20 sin( x3 ) dx 16
x 0 x 0 x 0 x 0
3 x 0 x 2
along BA : I 2 {sin( x ) dx 6 xydy} {sin( x3 ) 0}dx sin( x3 ) dx
BA x2 x 0
y 0
along AO : I 3 {sin( x 3 )dx 6 xydy} {sin(03 ) 0}dx 0
A0 y 2
x 2 x2
Hence I I1 I 2 I 3 sin( x 3 )dx 16 sin( x 3 ) dx 0 16
x 0 x 0
Example 2.279 [ME-2016 (2 marks)]: The value of the line integral C r ds , where C is a circle
F
of radius 4 units is _____. Here F ( x, y ) y iˆ 2 x ˆj and r is the UNIT tangent vector on the
curve C at an arc length s from a reference point on the curve. iˆ and ĵ are the basis vectors in the
x y Cartesian reference. In evaluating the line integral, the curve has to be traversed in the counter-
clockwise direction.
Solution: Any point on the curve C : x 2 y 2 a 2 is x a cos and y a sin with 0 2 ,
where a 4 . So we have r ( ) a cos iˆ a sin ˆj dr d a sin iˆ a cos ˆj . So the line
2
integral
F r ds F r ( ) r( )d {(a sin ) iˆ (2a cos ) ˆj} {a sin iˆ a cos ˆj}d
C C 0
2
2
C F r ds a 2 0 { sin 2 2 cos 2 }d a 2 0
2
1 3
cos 2 d a
2 2
2
3 sin 2
2 2 2 0
16
F r ds a 2 16 .
C
Example 2.280 [XE-2016 (2 marks)]: Calculate F dr F dr , where C1 : r (t ) (t , t 2 ) and
C1 C2
C2 : r (t ) (t , t ) , t varying from 0 to 1 and F xy ˆj .
Solution: For C1 , x t dx dt and y t 2 dy 2tdt ,
t 1
t 1 3 t 1 4 t5 2
ˆ ˆ ˆ ˆ
so F dr ( xy j ) ( dx i dy j dz k ) xy dy t (2tdt ) 2 t dt 2
C1 C1 C1 t 0 t 0
5 t 0 5
Path Independence of Line Integrals: The line integral given by Eq. 2.10 is
path independent in a domain D in space if every pair of endpoints A , B in D ,
the integral Eq. 2.10 has the same value for all paths in D that begin at A and end
at B , as shown in Fig. 2.48, i.e. a line integral (Eq. 2.10) with continuous F1 , F2
and F3 in a domain D in space is path independent in D if and only if
F F1 i F2 j F3 k is the gradient of some function f in D , i.e., Figure 2.48: Path
Independence
F grad f F1 f x , F2 f y , F3 f z and in this case
C ( F1dx F2 dy F3 dz) f ( B) f ( A)
Vector Fields that are Gradients of Scalar Fields (Potentials): Some
vector fields can be obtained from scalar fields. Such a vector field is
given by a vector function F( P) , which is obtained as the gradient of a
scalar function, F ( P) grad f ( P ) . The function f ( P ) is called a
potential function or a potential of F( P) . Such a F( P) and the
corresponding vector field are called conservative because in such a vector
field, energy is conserved; i.e. no energy is lost (or gained) in displacing a
Figure 2.49: Closed Path
body from a point P to another point in the field and back to P .
Path Independence and Integration around Closed Curves: The integral of Eq. 2.10 is path
independent in a domain D if and only if its value around every closed path in D is zero.
Path Independence and Exactness of Differential Forms: The path independence to the
exactness of the differential form given by
F dr F1 dx F2 dy F3 dz (2.11)
under the integral sign (2.10) is called exact in a domain D in space if it is the differential
df (f x) dx (f y )dy (f z )dz (grad f ) dr of a differentiable function f ( x, y, z )
everywhere in D .
Let F1 , F2 , F3 in the line integral curl F curl (grad f ) 0
f 0
C F(r ) dr C ( F1 dx F2 dy F3 dz) be continuous first partial i j k
derivatives in a domain D in space, then, if the differential form
x y z 0
F dr F1 dx F2 dy F3 dz is exact in D , and Eq. 2.10 is path
independent, then in D F1 F2 F3
F3 y F2 z , F1 z F3 x , F2 x F1 y (2.12)
Hence, if Eq. 2.12 holds in D and D is simply connected, then Eq. 2.11 is exact in D and
thus Eq. 2.10 is path independent.
Example 2.281 [ME-2005 (2 marks)]: The line integral
V d r of the vector function
V ( r ) 2 xyz iˆ x z ˆj x y kˆ from the origin to the point P(1,1,1) is
2 2
(a) 1 (b) 0 (c) –1 (d) cannot be determined without specifying the path
Solution (a): The given function is V ( r ) 2 xyz iˆ x z ˆj x y kˆ V iˆ V ˆj V kˆ .
2 2
1 2 3 As
V3 y V2 z , V1 z V3 x , V2 x V1 y the given function is exact and so the
line integral V d r is path independent. Now we have to find the potential function ( f ) of V such
that V grad f . Since, f x V1 f 2 xyz dx x 2 yz c , where c h( y , z ) is a function of y
and z . Now f y V2 x 2 z (c y ) x 2 z c y 0 c is a function of z only, i.e.,
c h( z ) . Now f z V3 x 2 y (c z ) x 2 y (c z ) 0 c is also not a function of z ,
2
thus c is just a constant and not a function of x or y or z ; so f x yz c ; if A(0, 0, 0) and
B(1,1,1) are endpoints of the given line integral V d r f ( B ) f ( A) (12 1 1 c) (0 c) 1
C
[Similar question was also asked in PI-2011 (2 marks)]
Example 2.282 [EC-2008 (2 marks)]: Consider points P and Q in (a) –1 (b) 0 (c) 1
the x y plane, with P(1, 0) and Q (0,1) . The line integral (d) depends on the direction
Q (clockwise or anti-
2 ( xdx ydy ) along the semi-circle with the line segment PQ as clockwise) of the
P
its diameter is semicircle
Q Q
Solution (b): The given integral can be written as I (2 xdx 2 ydy ) F (r ) dr , where
P P
Example 2.284 [EE-2016 (2 marks)]: The line integral of the vector field
F 5 xz iˆ (3 x 2 y ) ˆj x z kˆ along a path from (0, 0, 0) to (1,1,1) parameterized by (t , t , t ) is __.
2 2 2
xi xi xi 1 . This integral is represented by the area between the curve y f ( x) , the x axis and
between x a and x b . Now consider x f ( x , y ) and a region R of the xy plane as shown in
n
Fig. 2.50. The integral of f ( x, y ) over region R is R f ( x, y )dA lim
n ,all A 0
i
i 1 f ( xi , yi )Ai ,
where Ai ( i 1, 2, , n ) is a partition of R into n elements of area Ai and ( xi , yi ) is a point in
Ai . Now z f ( x, y ) represents a surface, and so f ( xi , yi )Ai zi Ai is the volume between z 0
and z zi on the base Ai . The integral R f ( x, y)dA is the limit of the sum of all such volumes,
and so it is the volume under the surface z f ( x, y ) above the region R . The partition of R into
elementary areas can be achieved using grid lines parallel to the x and y axes as shown in Fig. 2.51
(a). Then Ai xi yi and we can write
n
R f ( x, y )dA R f ( x, y )dx dy nlim
i1 f ( xi , yi ) xi yi (2.13)
Other partitions may be chosen, for example a polar grid as shown in Fig. 2.51 (b). Then the
element of area is (ri i )ri Ai and we write,
Figure 2.51: Grid for the partition of R (a) Rectangular Cartesian (b) Polar
Figure 2.50: Volume as an integral
We can evaluate integrals of the type R f ( x, y)dx dy as a repeated single integrals in x and y .
Consequently they are usually
called double integrals. Consider
the region R shown in Fig. 2.52 (a)
with boundary ACBD. Let the curve
ACB be given by y g1 ( x) and the
curve ADB by y g2 ( x) , then we
can evaluate R f ( x, y)dx dy by
summing for y first over the yi , Figure 2.52: Evaluation of double integrals for a Region R
holding x constant, from y g1 ( xi ) to y g2 ( xi ) and then summing all such strips from A to B ,
i.e. from x a to x b . Thus
Here the integral inside the brackets is evaluated first, integrating w.r.t. y , keeping the value of x
fixed, and then the result of this integration is integrated w.r.t. x .
Alternatively, as shown in Fig. 2.52 (b), we can sum for x first and then y . If the curve CAD is
represented by x h1 ( y ) and the curve CBD by x h2 ( y ) , thus
In both Eqs. 2.15 and 2.16, n min(n1 , n2 ) . If the double integral exists then these two results are
equal, and in going from one to the other we have changed the order of integration. Notice that the
limits of integration are also changed in the process.
1 1x x
Example 2.285 [CS-1993, EC-1993 (1 mark)]: Find value of the double integral dx dy .
0 0 1 y2
1 1x x 1 1x x 1 1x 1
Solution: I
0 0 1 y 2
dx dy
0 0 1 y 2 0
dy dx x tan 1 y 0 0
dx x tan 1 (1 x ) 0 dx
1 1
1 1 2 2 2
x dx (cot x )(1 2) x {1 (1 x )}(1 2) x dx
I x cot
0 0
II I
Example 2.286 [ME-2005 (1 mark)]: Changing the order of integration in the double integral
8 2 s q
I x 4 f ( x, y )dy dx leads to I p f ( x, y )dx dy . Then what is q ?
0 r
2
(a) 4y (b) 16y (c) x (d) 8
Solution (a): For the given integral, we integrate first w.r.t. y in
which the vertical length of element dx , as shown in figure, vary
from y 0 to y x 4 ; and then integrate w.r.t. x in which this
element dx shifts from x 0 to x 8 . Now, if we integrate first
w.r.t. x in which the horizontal length of element dy , as shown in
figure, vary from x 0 to x 4 y ; and then integrate w.r.t. y in
which this element dy shifts from y0 to y 2, thus
y 2 x4 y s q
I x0 f ( x, y )dx dy p f ( x, y )dx dy q 4 y
y0 r
So I xy dx dy
P
y 1
y0 x 2 2 y
x 0
xy dx dy
y 1
y0
(1 2)( x 2 y) xx 202 y dy
y 1 y 1 y 1
I
y 0
(1 2)(2 2 y ) 2
y 0 dy y 0
2(1 y ) y dy 2
2
y 0
(1 y 2 2 y ) y dy
y 1 y 1
I 2
y 0
( y y 3 2 y 2 ) dy 2 (1 2) y 2 (1 4) y 4 (2 3) y 3 y 0
(1 2) (1 4) (2 3) 0 1 6
[Similar question was also asked in ME-1997 (2 marks)]
3 x
Example 2.288 [CE-2008 (2 marks)]: The value of 0 0 (6 x y)dx dy is
(a) 13.5 (b) 27.0 (c) 40.5 (d) 54.0
2 2 2
Example 2.289: Evaluate R x ydA , where R is the region x y 1 .
Solution: As the domain of integration is a circle so, x r cos and y r sin and dA r d dr
I x 2 ydA
R
r 1
r 0 2
0
r 2 cos 2 r sin r d dr r 1
r 0 2
0
r 4 cos 2 sin d dr .
Note that in this integral the integration is such that we can separate the variables r and and hence
we can write, I
r 0
r 1
2
0
cos 2 sin d dr . As the limit of integration with respect to do not
r 0
r dr
0
2
cos2 sin d {(r 5 )10 } (1 3) cos3
2
0
0.
Example 2.290 [EE-2009 (2 marks)]: f ( x, y ) is a continuous function
2
defined over ( x, y ) 0,1 0,1 . Given the two constraints, x y and
y x2 , the volume under f ( x, y ) is
y 1 x y y 1 x 1
(a) y 0 x y 2
f ( x, y ) dx dy (b) y x x y
2 2
f ( x, y )dx dy
y 1 x 1 y x x y
(c) y0 x0 f ( x, y)dx dy (d) y 0 x0 f ( x, y) dx dy
2 2
Solution (a): The shaded region in the figure satisfy the two constraints x y and y x for all
y 1 x y
( x, y ) 0,1 0,1 , as shown in figure. So the volume under f ( x, y ) is V
y 0 x y 2
f ( x, y ) dx dy
1 x 2
Example 2.291 [XE-2013 (1 mark)]: The integral 0 x 2
( x y )e x y dy dx equals
(a) (e 2) e (b) (e 1) 2 e (c) (e 1) 2 (d) (e 2) 2e
Solution (d): As the given integral is first integrated w.r.t. y with y
2
changing from y x to y x ; and then w.r.t. x with x changing from
x 0 to x 1 . Now for changing the order of integration, i.e., if we
integrate first w.r.t. x in which the horizontal length of element dy , as
shown in figure, vary from x y to x y1 2 ; and then integrate w.r.t. y in
which this element dy shifts from y 0 to y 1 , thus
I
1
0 x
x x y
2 ( x y )e dy dx
2 y 1
y 0 x y
x y
( x y)e x y dx
2
y 0
y 1
( 1 2)e x
2
y
x y
x y
dy
y 1
I (1 2)
y 0
e 1
e y dy ( 1 2)(e 1 y e y ) yy 10 ( 1 2){(2 e) 1} (e 2) 2e
Example 2.292 [MA-2014 (2 marks)]: The value of R xy dx dy , where R is the region in the first
2
quadrant bounded by the curves y x , y x 2 and x 0 is …………….
Solution: The bounding curves intersect where x 2 x 2 , which gives x 1 (with y 1 ) and
x 2 (with y 4 ) in which A(1,1) lies in the first quadrant. So the region lies in the first quadrant
is AOBA with B (0, 2) and O (0, 0) . In this question, we choose to take y first because the formula
for the boundary is easier to deal with y x2 rather that x y1 2 . Thus we get,
I xy dx dy
R
x 1
x 0 y x
y 2 x
2
xy dy dx
x 1
x 0
(1 2) xy 2 y 2 x
y x 2
dx (1 2)
x 1
x 0
x(2 x) 2
xx 4 dx
x 1 x 1
I (1 2)
x 0
x 5
x 3 4 x 2 4 x dx (1 2) (1 6) x 6 (1 4) x 4 (4 3) x 3 (4 2) x 2 x0
0.375
2 x x y
Example 2.293 [ME-2014 (2 marks)]: The value of the integral 0 0 e dy dx is
2 2 2 2
(a) (e 1) 2 (b) (e 1) 2 (c) (e e) 2 (d) (1 2) e (1 e)
0 0 e
2 x 2 2 x2
Solution (b): I x y
dy dx 0 (e
x y y x
) y 0 dx 0 (e 2 x e x ) dx (1 2)e 2 x e x x 0
Example 2.294 [EC-2016 (2 marks)] The integral (1 2 ) ( x y 10) dx dy , where D denotes the
D
2 2
disc: x y 4 , evaluates to _____.
Solution: D : x 2 y 2 22 , converting the region D in polar coordinates, we have x r cos ,
y r sin , where r varies from 0 to 2; and varies from 0 to 2 . Now
2 2
I (1 2 ) ( x y 10)dx dy (1 2 )
r 0 0
( r cos r sin 10) rdrd
D
2 2 2 2
I (1 2 ) (r 2 cos r 2 sin 10r )d dr (1 2 ) (r 2 sin r 2 cos 10 r ) dr
r 0 0 r 0 0
2 2
I (1 2 ) 20 r dr (1 2 )20 rdr 10(1 2)(r 2 )02 20
r 0 r 0
Now, we can evaluate G 2 using the polar coordinates, whose differential area can be written as
r dr d , with limits r 0 to r and 0 to 2 . Thus
2 r2 2 2 2
d e r r dr 2 e r r dr . Now by letting r k 2rdr dk
2
G2 0 e r dr d
0 0 0 0
with limit changing from k0 to k ,
k
G 2 e k dk (1 ) (e k ) G . Thus
0 k 0
x2
e dx (2.17)
For a general quadratic exponent we simply complete the square and then integrate using a similar
2
ax 2 bx c a x (b 2a ) 2
4a ) c 2
4 a )c
change of variables, i.e., e dx e e (b dx ( a )e( b
a a x 2
Now, using definite integral property a f ( x ) dx 2 f ( x) dx if f ( x ) is an even function; as e
0
x2 2
is an even function thus from Eq. 2.14, we have e dx 2 e x dx . Thus
0
x2
0 e dx (1 2) (2.18)
3
Example 2.295 [ME-1994 (1 mark)]: Evaluate 0 ye y dy .
3 2 2 2 t 2 t
Solution: Let y t 3 y dy 2tdt dy 2
dt 43
dt ; also at y 0 , t 0 and at y ,
3y 3t
3 2 2 t 2 t 2 21
t . Thus, I ye y dy t1 3e t
3 0
43
dt e dt .
0 0 3t 32 3
Example 2.296 [EC-2005 (2 marks)]: The Solution (a): From Eq. 2.18 with 1 8,
value of the integral 1
I (1 2 ) exp x 8 dx is
2
0
exp x 2 8 dx
2
(1 8) 2
0
(a) 1 (b) (c) 2 (d) 2 I (1 2 ) 2 1
[Similar question was also asked in PI-2010 (1 mark)]
x2 y2
Example 2.297 [IN-2007 (2 marks)]: The value of the integral 0 0 e e dx dy is
0
x
2
0
y
2
2
1
Example 2.298 [XE-2013 (1 mark)]: The value of the integral 0 1 ( log e t ) dt is
(a) 2 (b) (c) (d) 2
2 2
2 k
Solution (b): Let log e t k (1 t ) dt 2kdk dt 2 k t dk 2ke dk ( t e k ); also at
0 2 2 1
t 0 , k 2 k and at t 1 , k 2 0 k 0 . I 2 e k dk 2 e k dk 2
0 2
Here we integrate along the entire boundary C of R in such a sense that Figure 2.53: Green's Theorem
R is on the left as we advance in the direction of integration, as shown in in a Plane
Fig. 2.53
Proof: Consider the first on the left side of the integral,
(F2 x) dx dy
R
yd
y c x g2 ( y )
x g1 ( y )
( F2 x) dx dy yd
y c
F2 g 2 ( y ), y F2 g1 ( y ), y dy
(F2 x ) dx dy F2 ( x, y )dy F2 ( x, y ) dy F2 ( x, y ) dy
F2 ( x, y )dy
R LMN LKN LMNKL C
Similarly, (F1 x ) dx dy
F1 ( x, y )dy .
R C
It follows from the Green’s theorem (Eq. 2.19) that the area A enclosed by the closed curve is
1
given by A 1 dx dy
x dy
y dy
2 C
( y dx x dy)
R C C
u v 2 v u v 2v u v u v
Eq. 2.19, A u u dx dy R dx dy .
R
x y xv y x yx x y y x
This implies that the element of area ( du dv ) is equivalent to the element
u v u v
dx dy . The modulus sign is used to preserve the orientation of the curve.
x y y x
x y x y . ( x, y )
Similarly ‘ dx dy J ( x, y ) du dv ,’ where J ( x , y )
(u , v ) u v v u
Hence, this enables us to make a general change of coordinates in a double integral as,
R f ( x, y)dx dy R f x(u, v), y (u, v) J ( x, y) du dv (2.20)
The Green’s theorem (Eq. 2.19) in a plane may be generalised to three
dimensions as A curl ( P, Q, 0) k dx dy C ( P, Q, 0) dr .
For a general surface S with bounding curve C , as shown in Fig.
2.54, this identity becomes S curl F(r ) dS C F(r ) dr , where
Figure 2.54: Three-
dS nˆ dS is the vector element of surface area and n̂ is a unit vector dimensional
along the normal. This generalisation is called Stokes’ theorem, and generalization of Green’s
will be discussed later, after introducing the concept of surface theorem
integral.
C ( xydy y dx) , where C
2
Example 2.299 [CE-2005 (2 marks)]: Value of the integral is the square
cut from the first quadrant by the line x 1 and y 1 will be (use Green’s theorem to change the line
integral into double integral)
(a) 1 2 (b) 1 (c) 3 2 (d) 5 3
Solution (c): Here F1 ( x, y ) y 2 and F2 ( x, y) xy , so that F1 y 2 y and F2 x y
F2 x F1 y y (2 y ) 3 y . Thus the line integral transforms into an easy double integral
by using Green’s theorem in a plane, I ( xydy y dx) (F2 x F1 y ) dx dy
2
C R
y 1 x 1 y 1 y 1 y 1
I 3 y dx dy (3 yx ) xx 10 dy 3 y dy (3 2) y 2 3 2
y 0 x 0 y0 y 0 y 0
y 2
8 ( y 2) 1
Example 2.302 [EE-2014 (2 marks)]: To evaluate the double integral (1 2)(2 x y ) dx dy
0
J
x u
y u
x v
y v
1
0
1
2
2;
8
0 ( y 2) 1
y 2
2x y
2
dx dy
4
0 u J du dv 2u du dv
1
0 0
4 1
Solution (c): As given in Eq. 2.19, we integrate along the entire boundary C of R in such a sense
that R is on the left as we advance in the direction of integration; this criteria is shown by the figure
given in option (c).
Example 2.305 [EC-2016 (2 marks)]: Suppose C is the closed curve defined as the circle
x 2 y 2 1 with C oriented anti-clockwise. The value ( xy dx x
2 2
ydy ) over the curve C equals _.
F F
Solution: Using Green’s theorem in a plane, which is C ( F1dx F2 dy) R x2 y1 dxdy , where
F1 and F2 are functions of ( x, y ) . So we have F1 xy 2 and F2 x 2 y , thus
( x 2 y ) ( xy 2 )
C ( xy dx x ydy ) R
2 2
dxdy R 2 xy 2 xy dxdy 0 .
x y
Parametric Coordinated of Sphere: When evaluating surface integral over the surface of a
sphere of radius a , as shown in Fig 2.56 (a), the parametric coordinates can be given as,
x a sin cos , y a sin sin , z a cos and dS a 2 sin d d .
Parametric Coordinates of a Cylinder: When evaluating surface integral over the surface of a
cylinder of radius
a , as shown in
Fig. 2.56 (b), the
parametric
coordinates can be
given as,
x a cos ,
y a sin , z z
, dS a dz d .
A surface S is
called a smooth
surface if its
surface normal
depends Figure 2.56: (a) Surface element in spherical polar coordinates (b) Surface element in
continuously on cylindrical polar coordinates
the points of S . S is called piecewise smooth if it consists of finitely many smooth portions.
r (t ) r u (t ), v (t ) . By differentiation and the use of the chain rule, we obtain a tangent vector of C
dr dr dr dr dr
on S , as r (t ) u v . Hence the partial derivatives ru at P are and rv
dt du dv du dv
tangential to S at P . We assume that they are linearly independent, which geometrically means that
the curves u const and v const on S intersect at P at a non-zero angle. Then ru and rv span the
tangent plane of S at P . Hence their cross product gives a normal vector N of S at P , as
N ru rv 0 (2.22)
The corresponding unit normal vector n (as shown in Fig. 2.57) of S at P is given as,
n 1 N N 1 ru rv (ru rv ) (2.23)
Also, if S is represented by g ( x, y , z ) 0 , then n can also be given as,
n 1 grad g (grad g ) (2.24)
Surface Integrals of a vector function: To define a surface integral, we take a surface S , given
by a parametric representation as given by Eq. 2.21. We also assume S to be piecewise smooth so
that S has a normal vector given by Eq. 2.22 and unit normal vector given by Eq. 2.23. For a given
vector function F , we can now define the surface integral over S by
S F ndA R F r(u, v) N(u, v) du dv (2.25)
Here N N n and N ru rv is the area of the parallelogram with sides ru and rv , by the
definition of cross product; and ndA n N du dv N du dv dA N du dv is the element of area of
S . Also, F n is the normal component of F . In component form, F F1 i F2 j F3 k ,
N N1 i N 2 j N 3 k and n cos i cos j cos k , ( , , are the angles between n and
the coordinate axes x , y , z axes, respectively), we can write Eq. 2.25 as,
S F n dS A ( F1 cos F2 cos F3 cos )dA R ( F1 N1 F2 N2 F3 N3 )du dv . We can write
cos dA dy dz , cos dA dz dx and cos dA dx dy , then above integral becomes,
S F n dS S ( F1 dy dz F2 dz dx F3 dx dy ) (2.26)
We can use Eq. 2.26 to evaluate surface integrals by converting them to double integrals over regions
in the coordinate planes of the xyz coordinate system. But we must carefully take into account the
orientation of S (the choice of n ). For e.g. if the surface S is given by z h( x, y ) with ( x , y )
varying in a region R in the xy plane and if S is oriented so that cos 0 then,
S F3 cos dA S F3 dx dy R F3 x, y, h( x, y) dx dy ; but if cos 0 , the integral on the right
side gets a minus sign.
The surface integral arises naturally in flow problems, where it gives the flux across S ( mass
of fluid crossing S per unit time) when F v , where is the density of the fluid and v is the
velocity vector of the flow. We may thus call the surface integral as the flux integral.
is the triangle bounded by the two coordinate axes and the straight line x y 1 , obtained from
x y z 1 by setting z 0. Thus 0 x 1 y , 0 y 1. Now ru dr du i k ,
rv dr dv j k N ru rv (i k ) ( j k ) i j k
F N (u 2 i 3v 2 k ) (i j k ) u 2 3v 2 . So,
v 1 u 1 v v 1
S F n dA R (u
2
3v 2 )dudv
v 0 u 0
(u 2 3v 2 ) du
v 0
(1 3)(1 v) 3
3v 2 (1 v) dv 1 3
Here R * is the projection of S into the xy plane, as shown in Fig. 2.58 and the normal vector N
on S points up. If it points down, the integral on the right is preceded by a minus sign. Eq. 2.27 with
G (r ) 1 gives the area A( S ) of S : z f ( x, y ) and thus,
A( S ) 1 dA 1 (f x ) 2 (f y ) 2 dx dy (2.28)
S R*
where R * is the projection of S into the xy plane.
Example 2.307: Evaluate the surface integral S ( x y z )dS , where S is the portion of the sphere
x 2 y 2 z 2 1 that lies in the first quadrant.
Solution: The surface S is illustrated in the Figure (a). Taking z 1 x 2 y 2 , we have
2 2 2 2
z x x 1 x 2 y 2 and z y y 1 x 2 y 2 1 ( zx ) ( z y ) 1 (1 x y )
x 1 x2 y 2 y
1 x 2 y 2 1 dy dx
2
1 x
0
1
I x sin 1 y
1 x 2 1 x2 y 2 y
0
1
dx ( 2) x 2 1 x 2 dx
0
1
I ( 4) x 2 x 1 x 2 sin 1 x 3 4
0
2 2
I ( x y z ) dS I
S 0 0 sin cos sin sin cos sin d d
2 1
I (1 4) cos (1 4) sin 2 d 3 4
0
Example 2.308 [MN-2008 (2 marks)]: The direction of gradient vector at a point (1,1, 2) on a
surface S ( x, y , z ) x 2 y 2 z is
(a) (1 3)(2iˆ 2 ˆj kˆ) (b) (1 3)( 2iˆ 2 ˆj kˆ)
(c) (1 3)(2iˆ 2 ˆj kˆ ) (d) (1 3)(2iˆ 2 ˆj kˆ )
Solution (d): The direction of gradient vector is the outward unit normal vector of the surface at a
point (1,1, 2) on a surface S ( x, y , z ) x 2 y 2 z . n S (S x)i (S y ) j (S z )k
n S 2 x i 2 y j z k n P 1,1,2 2 i 2 j 1 k nˆ (1 3)(2 i 2 j 1k ) .
Example 2.309 [ME-2012, PI-2012 (1 mark)]: For the spherical surface, x 2 y 2 z 2 1 , the unit
outward normal vector at the point P (1 2 ,1 2 , 0) is given by
(a) (1 2)iˆ (1 2) ˆj (b) (1 2)iˆ (1 2) ˆj (c) k̂ (d) (1 3)iˆ (1 3) ˆj (1 3) kˆ
Solution (a): For the surface f ( x, y , z ) x 2 y 2 z 2 1 , the outward normal vector
n f (f x)i (f y ) j (f z )k 2 x i 2 y j 2 z k n P 1 2,1 2,0 2i 2 j.
So the unit outward normal vector at P (1 2 ,1 2 , 0) is nˆ n n (1 2) i (1 2) j
[Similar questions were also asked in CH-2008, IN-2009 (1 mark), TF-2008 (2 marks)]
Example 2.310 [XE-2016 (1 mark)]: The value of the surface integral F n dS over the sphere
2 2 2 ˆ ˆ
given by x y z 1 , where F 4 x i z k , and n denotes the outward unit normal, is
(a) (b) 2 (c) 3 (d) 4
Solution (d): Using divergence theorem, div F dV F dS F n dA , we have
V S S
ˆ ˆ
F 4 x i z k , so div F F ( x )(4 x) ( y )(0) ( z )( z ) 4 1 3 . Now
3
S F n dA V div F dV V 3 dV 3V 3(4 3) r 4 (since r 1 ). So option (d) is correct.
Figure 2.61: (a) Partition of region T into volume elements ∆ (b) The volume integral in terms of rectangular Cartesian
coordinates.
Note that there are six different orders in which the integration in Eq. 2.35 can be carried out. As we
saw for double integrals, the expression for the element of volume dV dx dy dz under the
transformation x x (u , v, w) , y y (u, v, w) and z z (u , v, w) may obtained using the Jacobian as,
x u y u z u
( x, y , z )
J x v y v z v , so dV dx dy dz J du dv dw (2.36)
(u , v, w)
x w y w z w
In case of cylindrical polar coordinates ( r , , z ) , as shown in Fig. 2.62 (a), x r cos ,
cos sin 0
y r sin and z z J r sin cos 0 r dV r dr d dz
0 0 1
In case of spherical polar coordinates ( r , , ) , as shown in Fig. 2.62 (b), x r sin cos ,
sin cos sin sin cos
y r sin sin and z r cos J r cos cos r cos sin r sin r 2 sin
r sin sin r sin cos 0
dV r 2 sin dr d d
Figure 2.62: Volume element in (a) cylindrical (b) spherical polar coordinates
Example 2.311 [ME-2004 (2 marks)]: The volume of an object expressed in spherical co-ordinates
2 3 1 2
is given by V 0 0 r sin dr d d . The value of the integral is
0
Example 2.312 [EC-2014 (2 marks)]: The volume under the surface z ( x , y ) x y and above the
triangle in the x y plane defined by { 0 y x and 0 x 12 } is …………….
Solution: Taking elemental strip of length dx whose height vary from z 0 to z x y ; as this
strip lies in the triangle in xy plane which is defined by 0 y x and 0 x 12 ; so its width is
changing from y 0 to y x ; and whose length is changing from x 0 to x 12 . So
V
x 12
x 0
y x
y 0
z x y
z 0
dz dy dx
x 12
x 0 y x
y 0
z zz 0x y dy dx x 12
x 0
yx
y0
( x y ) dy dx . Thus
X x Y x Z x 0 2 1
( X , Y , Z )
J X y Y y Z y 3 0 3 0 2(0 12) 1(9 0) 15
( x, y , z )
X z Y z Z z 4 3 0
(2 X Y 2 Z ) dX dY dZ 5 z 15 dx dy dz 75 z dx dy dz 15
g(S) S S
Example 2.315 [EC-2016 (2 marks)]: A triangle in the xy plane is bounded by the straight lines
2 x 3 y , y 0 and x 3 . The volume above the triangle and under the plane x y z 6 is _____.
x 0
y2 x 3
y 0 z dx dy
x 3
x 0 y 2 x 3
y0
(6 x y ) dy dx
y 2 x 3 3
x 3 y2
x 3 8x2 2 8x3 2 8(3)3
V 6 y xy dx 4 x dx 2 x 2(3) 10 .
x 0
2 y 0 x 0
9 27 0 27
F F2 F3
V x1
y
dx dy dz S ( F1 dy dz F2 dz dx F3 dx dy )
z
(2.40)
Common Data for Questions 2.317 & 2.318: Consider the vector field
A ( y z )iˆ ( x z ) ˆj ( x y ) kˆ defined over the unit sphere x y z 1 .
3 3 3 3 3 3 2 2 2
Example 2.317 AE-2009 (4 marks): The surface integral (taken over the unit sphere) of the
component A normal to the surface is
(a) (b) 1 (c) 0 (d) 4
Solution (c): As div A A ( x)( y 3 z 3 ) ( y )( x 3 z 3 ) ( z )( x 3 y 3 ) 0 , so by using
divergence theorem S A ndA div A dV 0dV 0 .
V V
Example 2.318 AE-2009 (4 marks): The magnitude of the component of A normal to the spherical
surface at the point 1 3 ,1 3 ,1
3 is
(a) 1 3 (b) 2 3 (c) 3 3 (d) 4 3
Solution (b): As normal to the given spherical surface, f ( x, y , z ) x 2 y 2 z 2 1 , is
n f ( x)i ( y ) j ( z )k 2 x i 2 y j 2 z k n 1 3,1 3,1 3 (2 3)( i j k ) ; also at
the given point, A {2 (3 3 ) (i j k ) . As n A ( 2 3 )(i j k ) {2 (3 3 )}(i j k ) 4 3 and
n (2 3) 12 12 12 2 . Now, the component of A normal to the spherical surface is same as
n A nA nA 43 2
the component of A along n which is 2
n , whose magnitude is 2
n .
n n n 2 3
Example 2.319 [IN-2014 (1 mark)]: A vector is defined as f y iˆ x ˆj z kˆ where iˆ , ĵ and k̂ are
unit vectors in Cartesian ( x, y , z ) coordinate system. The surface integral
f dS over the closed
surface S of a cube with vertices having the following coordinates: (0, 0, 0) , (1, 0, 0) , (1,1, 0) ,
(0,1, 0) , (0, 0,1) , (1, 0,1) , (1,1,1) , (0,1,1) is …………….
Solution: div f f ( x )i ( y ) j ( z )k y i x j z k 1
So using divergence theorem,
S f dS
V div f dV
V 1dV
V 1dV volume of the
given cube which is 1 (as side of cube is of unit length). So answer is 1.
Example 2.320 [XE-2011 (1 mark)]: A vector field is called solenoidal if its divergence is zero.
Consider the vector fields P and Q given by P ( x, y , z ) k iˆ k ˆj k kˆ , where k (2 x 2 8 xy 2 z )1 2 3 1
3 2 2 3
2
; k2 (3 x y 3 xy ) ; k3 (4 y z 2 x z ) and Q ( x, y , z ) xyz P ( x, y , z ) . Then
(a) P and Q are both solenoidal (b) Both P and Q are not solenoidal
(c) P is solenoidal but not Q (d) Q is solenoidal but not P
Solution (d): div P ( x )(2 x 2 8 xy 2 z ) ( y )(3 x 3 y 3xy ) ( z )( 4 y 2 z 2 2 x 3 z )
div P 4 x 8 y 2 z 3 x 3 3 x 8 y 2 z 2 x 3 x x 3 P is not solenoidal.
divQ ( x) xyz 2 (2 x 2 8 xy 2 z ) ( y ) xyz 2 (3x 3 y 3 xy ) ( z ) xyz 2 (4 y 2 z 2 2 x 3 z )
3 2 2 3 3 4 2 2 2 2 2 3 4 4 3
div Q ( x )(2 x yz 8 x y z ) ( y )(3 x y z 3x y z ) ( z )( 4 xy z 2 x yz )
2 2 3 3 4 2 2 2 3 3 4 2
div Q 6 x yz 16 xy z 6 x yz 6 x yz 16 xy z 6 x yz 0 Q is solenoidal.
S e iˆ 3 y ˆj ze x kˆ ndA
x
Example 2.321 [XE-2012 (2 marks)]: Evaluation of ˆ over a surface
2 2 2
S : x y z 1 using Gauss divergence theorem gives
(a) 0 (b) 4 (c) 4 3 (d) 12
F e iˆ 3 y ˆj ze kˆ
x x
Solution (b): The given spherical surface has radius r 1 . For
div F F ( x)iˆ ( y ) ˆj ( z ) kˆ e x iˆ 3 y ˆj ze x kˆ 3 . Thus using divergence
3
theorem, ˆ F dS div F dV 3dV 3V 3(4 3) r
S F ndA 4
S V V
[Similar questions were also asked in CH-2009, AE-2013 (2 marks)]
Example 2.322 [ME-2013, PI-2013 (2 marks)]: The following surface integral is to be (a)
evaluated over a sphere for the given steady velocity vector field, F x iˆ y ˆj z kˆ (b) 2
(c) 3 4
defined with respect to a Cartesian coordinate system having iˆ , ĵ and k̂ as unit base
(d) 4
vectors. (1 4) F n dA , where S is the sphere, x 2 y 2 z 2 1 and n is the outward
S
unit normal vector to the sphere. The value of the surface integral is
Solution (a): div F F ( x )i ( y ) j ( z )k x i y j z k 3 . Using divergence
1 1 1 1 3 34 3
theorem, S 4 F ndA 4 S F dS 4 V div F dV 4 V 3dV 4 V 4 3 r 4 (Since r 1
for the given sphere)
Furthermore r dr ds is the unit tangent vector and s the Figure 2.64: Direction of and orientation
arc length of C . In component form, Eq. 2.43 can be written of
as,
F3 F F F3 F F
R y z2 N1 z1 x N 2 x2 y1 N3 du dv C ( F1dx F2 dy F3dz ) (2.43)
2
1. If C is a line segment from (3, 4, 0) to (1, 4, 2) , then C ( z y ) ds _____.
2. If C is the part of the curve x cos y from (1, 2 ) to (1, 0) and F y i 2 x j , then
C y dx xdy , where
2 2
4. Evaluate the line integral C is the arc of the parabola x 4 y from
( 5, 3) to (0, 2) . _____.
C (3 y e
sin x
5. Evaluate ) dx (7 x y 4 1)dy , where C is the circle x 2 y 2 9 .
(a) 4 (b) 8 (c) 16 (d) 36
6. Evaluate C (3x 5 y )dx ( x 6 y)dy , where C is the ellipse (1 4) x 2 y 2 1 in the
anticlockwise direction.
(a) 4 (b) 8 (c) 12 (d) 16
2 2 2 2
7. Evaluate S z ds , where S is the hemisphere given by x y z 1 with z 0 .
(a) 3 (b) 2 3 (c) 3 3 (d) 4 3
8. The area of the ellipse cut on the plane 2 x 3 y 6 z 60 by the circular cylinder x 2 y 2 2 x .
(a) 7 6 (b) 5 6 (c) 3 6 (d) 6
4 2 2 2
9. Evaluate I ( x y y z xz )dS , where S is the entire surface of the sphere
S
x2 y 2 z2 1 .
(a) 2 15 (b) 4 15 (c) 6 15 (d) 8 15
10. Find I F ndS , where F 2 x i 2 y j k , n is the outward pointing unit normal and S is
S
the entire surface consisting of S1 the part of the paraboloid z 1 x 2 y 2 with z 0 together
with S 2 disc {( x , y ) : x 2 y 2 1} .
(a) (b) 2 (c) 3 (d) 4
11. Which of the following vector fields are (a) Only V
conservative? (b) Only W
V (2 x 3 y z ) i ( 3x y 4 z ) j (4 y z ) k and (c) Both V and W
W (2 x 4 y 5 z ) i ( 4 x 2 y ) j ( 5 x 6 z ) k (d) Neither V nor W
12. Find the potential function ( ) for F (2 x 4 y 5 z ) i ( 4 x 2 y ) j ( 5 x 6 z ) k .
(a) x 2 4 yz 5 zx y 2 x 3 z 2 (b) y 2 4 yz 5 zx y 2 x 3 z 2
(c) x 2 4 yz 5 zx y 2 x 3 y 2 (d) x 2 yz zx y 2 x 3 z 2
13. Evaluate C F dr , where C is the curve from A(1, 0, 0) to B (0, 0,1) in which the plane x z 1
(a) L
4
(b) L4 2 (c) L4 3 (d) L4 4
17. R ( x y 1)dxdy , where R is the region inside the unit square in which x y 0.5 .
(a) 7 5 (b) 5 7 (c) 8 7 (d) 7 8
1 1
18. Evaluate 0 0 x max( x, y )dydx .
(a) 3 8 (b) 8 3 (c) 8 7 (d) 7 8
19. The volume of the solid bounded above by the plane z 4 x y and below by the rectangle
R {( x, y ) : 0 x 1, 0 y 2} is _____.
20. Evaluate D (3 x y )dA , where D is the triangle in the ( x, y ) plane bounded by the x axis
and the lines y x and x 1 .
2
21. EvaluateD (4 x 2)dA , where D is the region enclosed by the curves y x and y 2 x .
3
22. Evaluate ( xy y ) dA , where D is the region consisting of the square
D
{( x, y ) : 1 x 0, 0 y 1} together with the triangle {( x , y ) : x y 1, 0 x 1} .
(a) 23 40 (b) 23 40 (c) 40 23 (d) 40 23
23. Evaluate D (sin x) x dA , where D is the triangle {( x, y) : 0 y x, 0 x } .
24. Find the volume of the tetrahedron that lies in the first octant and is bounded by the three
coordinate planes and the plane z 5 2 x y . _____
2
25. Evaluate D ( x y ) dxdy , where D is the parallelogram bounded by the lines x y 0 ,
x y 1 , 2 x y 0 and 2 x y 3 .
26. Let D be the region in the first quadrant bounded by the hyperbolas xy 1 , xy 9 and the lines
y x , y 4 x ; then evaluate D
y x xy dxdy .
(a) 8 (52 3) ln 2 (b) 8 (52 3) ln 3 (c) (52 3) ln 3 (d) (52 3) ln 2
1 ey
27. Change the order of the integral 0 1 f ( x, y )dxdy .
1 log x e log x
(a) 0 1 f ( x, y)dydx (b) 0 1 f ( x, y )dydx
e 1 e 1
(c) 1 log x f ( x, y)dydx (d) 0 log x f ( x, y)dydx
1 1 2
28. Evaluate 0 x e y dydx [Hint: evaluate by changing order of integration]
(a) (e 1) 2 (b) (e 1) (c) (e 1) 3 (d) (e 1) 4
5 2 1
29. Change the order of integration 2 sin x f ( x, y )dydx
1 sin 1 y 2 1 sin 1 y 2
(a) 1 sin 1
y
f ( x, y ) dxdy (b) 1 sin 1
y
f ( x, y ) dxdy
1 sin 1 y 2 1 sin 1 y 2
(c) 0 sin 1
y
f ( x, y ) dxdy (d) 0 sin 1
y
f ( x, y ) dxdy
y 8 . _____
C ( xydx x
2
34. Evaluate y 3dy ) where C is the triangle with vertices (0, 0) , (1, 0) , (1, 2) with
positive orientation. _____
C ( y dx x dy )
3 3
35. If C is the positively oriented circle of radius 2 centred at the origin then
C ( y dx x dy) , where
3 3
36. Evaluate C are the two circles of radius 2 and radius 1 centred at the
origin with positive orientation.
(a) 20 (b) 20 (c) 22.5 (d) 22.5
2
37. Evaluate S curl F dS , where F z i 3 xy j x y k and S is the part of z 5 x 2 y 2 above
3 3
2.8.1 Sequence
A sequence is a list of numbers written in a specific order. The list may or may not have an infinite
number of terms in them. General sequence terms are denoted as follows: a1 (first term), a2 (second
term), , an ( nth term), an 1 ( ( n 1)th term), . As we are dealing with infinite sequences, each
term in the sequence will be followed by another term as given above. There are many ways for
denoting a sequence; the sequence {a1 , a2 , , an , an 1 , } can also be denoted as {an } or {an }n1 , ;
while, the sequences denoted by {an } or {an }n1 , an is usually denoted by a formula. For e.g., the
sequence {an }n1 , an ( 1) n 1 2n is written as 1,1 2 , 1 4 ,1 8, 1 16 , .
If an an 1 for every n , then the sequence is increasing; on the other hand if an an 1 for every
n , then the sequence is decreasing.
If a number m s.t. m an n , then the sequence is bounded below; also if there exists a
number M such that an M for every n , we say that the sequence is bounded above.
If the sequence is both bounded below and bounded above we call the sequence bounded.
Example 2.325: Determine if the following sequences are monotonic (i.e., increasing or decreasing)
2 n 1
and/or bounded. (a) n n 0
(b) ( 1) n 1
2 2
Solution: (a) As n ( n 1) n , i.e., an an 1 , so the sequence is a decreasing sequence and
hence monotonic. Also, as the sequence terms will be either zero or negative, thus the sequence is
bounded above but not bounded below and hence the sequence is not bounded. (b) For the sequence
n 1
given by ( 1) n 1
, the sequence terms are alternate 1 or 1 so the sequence is neither increasing
nor decreasing so it is not a monotonic sequence. Also the given sequence has terms either 1 or 1
so the sequence is bounded below by 1 and bounded above by 1 .
Limit of a sequence
If lim an L , where L is a finite number, then the value of the an L as n .
n
If lim an , then it means that the value of the an gets larger and larger as n approaches to .
n
If lim an , then it means that the value of the an gets negative larger and larger as n .
n
If lim an does not exist or is infinite then the sequence is diverges. If lim an we will say the
n n
lim (1 n ) ln b
Solution (b): Let L lim n1 n ln L lim(1 n) ln n L e n . As lim (1 n) ln n , takes
n n n
Example 2.327: Determine if the following sequences converge or diverge. If the sequence converges
determine its limit. (a) (3n 1) (10n 5n )
2
(b) e n
2
n 2
2n
n 1
3n 1 2 n 3 (1 n )
2
3 (1 n ) 3 ,
2 2
Theorems on Sequences
Given the sequence {an } , if we have a function f ( x ) s.t. f ( n) an and lim f ( x ) L then
x
lim an L , i.e. we take the limits of sequences much like we take the limit of functions. Thus all
x
algebraic properties from the limit of functions will also hold for the limit of sequences, i.e.,
lim (an bn ) lim an lim bn lim can c lim an
n n n n n
n
lim (a n bn ) lim an
n
lim b
n
n lim an bn lim an
n
n
lim b ,
n
n
p
provided lim bn 0
p
lim( an ) lim an
n n
, provided an 0 n
The sequence {r n }n 0 converges if 1 r 1 and diverges for all other value of r . Also,
lim r n 0 for 1 r 1 ; and lim r n 1 for r 1 .
n n
If both lim a2 n L and lim a2 n 1 L then the sequence {an } is convergent and lim an L
n n n
2.8.2 Series
Series is the sum of all terms in a sequence. Let a sequence given by {an }i1 , then s1 a1 ,
n
s2 a1 a2 , , sn a1 a2 an i 1 ai , . The sn are called partial sums and they will
form a sequence {sn }n1 . The limit of the sequence of partial sums {sn }n1 is defined as
n
lim sn lim i 1 ai i 1 ai , where
n n
i1 ai is called an infinite series (note that the series starts at
i 1 as the original sequence {an }i1 also starts at i 1 ; if the original sequence starts at i 2 then
the series will also starts at i 2 ).
Example 2.328: Write n1n2
(1 3n1 ) as a series that starts at n 3 .
Solution: As we want to increase the initial value by 2 and so all the n ’s in the series term must
decrease by 2, i.e., n1n2
(1 3n1 ) n3 ( n 2) 2 (1 3( n 2)1 ) n3 ( n 2) 2 (1 3n1 ) .
Divergence test: If lim an 0 then an will diverge; but if lim an 0 then an converges
n n
or diverges.
Integral test: Suppose that f ( x ) is continuous, positive and decreasing function on the interval
[ k , ) and that f ( n) an then
If k f ( x)dx n k an is also convergent
is convergent then
is divergent then n k an is also divergent
If f ( x) dx
k
The −series test: If k 0 then n k (1 n p ) converges if p 1 and diverges if p 1 [This
point was asked in AG-2016 (1 mark)]
Comparison test: As it is not possible to integrate every function so integral test is not use-full
all the times. We have another test for comparing the convergence or divergence of a series called
comparison test or limit comparison test which states that: if we have two series an and bn
with an , bn 0 for all n and an bn for all n , then
If bn is convergent then an is convergent
If an is divergent then bn is divergent
Limit Comparison test: Suppose that we have two series an and bn with an 0 , bn 0
for all n . Let us define c lim (an bn ) , if c is positive and finite then either both series converge
n
or both series diverge.
Alternating series test: Suppose that we have a series an and either an ( 1) n bn or
an ( 1) n 1 bn where bn 0 for all n . Then if lim bn 0 and {bn } is a decreasing sequence then
n
an 1 (n 1)!2 (2n)! ( n 1) 2 (1 1 n)
L lim lim 2
lim lim lim( n 1) 1 . Thus the
n an n (2n 1)! ( n !) n (2n 1) n (2 1 n) n
given series diverges.
Example 2.331 [TF-2009 (2 marks)]: If S is the largest possible set of real numbers x for which the
series n1 x n n is convergent, then S is
(a) ( 1,1) (b) [ 1,1) (c) [ 1,1] (d) (, )
Solution (a): Using the ratio test, the series n1 an , where an x n n , converges if
an 1 x n 1 n x
L lim 1 lim n
1 lim 1 x 1 x ( 1,1)
n an n n 1 x n 1 1 n
Example 2.332 [XE-2009 (2 marks)]: The infinite series m1(1)m x2
(1 x 2 ) m is
(a) Divergent for all x (b) Convergent only for x 1
(c) Convergent for all x (d) Divergent only for 1 x 1
Solution (c): The given infinite series can be written as m1 am m1 (1)m bm , where
bm x2 (1 x 2 )m . As lim bm lim x 2 (1 x 2 ) m x 2 (1 x 2 ) 0 ; also x 2 1 x 2 (1 x 2 ) m
m m
for all x and m [1, ) thus x 2 (1 x 2 ) m x 2 (1 x 2 ) m 1 for all x and m [1, ) thus
bm x 2 (1 x 2 )m is a decreasing sequence thus the given sequence converges for all x .
Example 2.333 [IN-2011 (2 marks)]: The series m0 (1 4m )( x 1)2 m converges for
(a) 2 x 2 (b) 1 x 3 (c) 3 x 1 (d) x 3
Solution (b): Using the ratio test, the series n1 an , where an (1 4 )( x 1) 2 m , converges if
m
am1 ( x 1) 2( m1) 4m ( x 1) 2 ( x 1) 2
L lim 1 lim m 1
2m
1 lim 1 1
m am m 4 ( x 1) m 4 4
1 (1 4)( x 1) 2 1 . Now we have two inequalities (i) (1 4)( x 1) 2 1 x and
(ii) (1 4)( x 1) 2 1 ( x 1) 2 4 x 2 2 x 1 4 x 2 2 x 3 0 ( x 1)( x 3) 0 . Now
using wavy curve method ( x 1)( x 3) 0 x ( 1, 3)
Example 2.334 [MT-2013 (2 marks)]: Which one of the following series is divergent?
(a) n 11 3n1 (b) n1 (1 n) (c) n 0 (1 2n ) (d) n1 (1 n n )
Solution (b): For option (a): As infinite geometric series, which converges to a (1 r ) , where a is
the first term and r is the common ratio. For the given series The given series n 11 3n1 , a 1
and r 1 3 . Thus the given series converges to 1 (1 1 3) 3 2 . The option (b) diverges which is
described in Example 2.337. For option (c): Again we have infinite geometric series whose first
term is a 1 and common ratio r 1 2 . Thus the given series converges to 1 (1 1 2) 2 . For
option (d): We have an 1 n n , so using the ratio test,
n n n
an 1 n n 1 1
L lim lim n 1
lim n
lim lim 1 0 0 . So we
n an n (n 1) n ( n 1) ( n 1) n
(1 1 n) n n 1
have L 1 and hence the given series is convergent.
Example 2.335 [AG-2014 (2 marks)]: Which of the following statements is true for the series given
below? sn 1 (1 2) (1 3) (1 4) (1 n)
(a) sn converges to log( n ) (b) sn converges to n
(c) sn converges to exp( n ) (d) sn diverges
Solution (d): Using the p series test we have n k (1 n p ) , where k 1 and p 1 2 1 , thus the
given series diverges.
Example 2.336 [MT-2014 (2 marks)]: The following power series P ( x ) will converge absolutely if
x is less than: P ( x) n 0 4n x n
an 1 4n 1 x n 1 4
n 1
x
n 1
1
lim 1 lim n n
1 lim n n
1 lim 4 x 1 4 x 1 x
n an n 4 x n 4 x n 4
Value of a Series: Let we have a series n1 an which converges to s , i.e., if we take the partial
n
sums then sn i 1 ai and this will form a convergent sequence and its limit will be s , i.e.,
lim sn s . Now limit converges means that we can make the partial sums, sn , as close to s as we
n
want simply by taking n large enough. In other words, if we take n large enough, then sn s . This
is one method of estimating the value of a series. We can just take a partial sum and use that as an
estimation of the value of the series. Now there are two questions: (i) How good is the estimation? (ii)
Is there any way to make the estimate better? So, for determining how good the estimation is, consider
n
the full series and strip out the first n terms, i.e., i1 ai i1 ai in 1 ai . The first series is the
partial sum sn ; and the second series (the one starting at i n 1 ) is called the remainder and denoted
by Rn . Thus s sn Rn Rn s sn . So the remainder Rn tells us the error, between the exact
value of the series and the value of the partial sums that we are using as the estimation of the value of
the series. As we do not have the actual value of series so we cannot get actual value of the remainder.
However, we can use some of the tests that we have got for convergence to get a good estimate of the
remainder provided we make some assumptions about the series. Once we have got an estimate on the
value of the remainder we will also have an idea on just how good a job the partial sum does of
estimating the actual value of the series. There are several tests that will allow us to get estimates of
the remainder, which are described as:
Integral Test: In this case we will need to assume that the series terms are all positive and will
eventually be decreasing. In definite integral, using rectangle method, we derived the integral
b
a f ( x) dx is the area under the curve y f ( x) bounded by x axis, x a and x b . Similarly,
the series could be thought of as an estimation of the area under the curve of f ( x ) where
f ( n) an . We can do something similar with the remainder, as Rn i n 1 ai an 1 an 2 .
If we start at x n 1 , taking rectangles of width 1 and use the left endpoints as the height of the
rectangle we
can estimate
the area
under f ( x )
on the
interval
[ n 1, ) as
shown in
Fig. 2.66(a).
We can see Figure 2.66: Remainder for Integral Test
that the
remainder, Rn , is the area estimation and it will overestimate the exact area. So, Rn f ( x )dx .
n 1
Next, we could also estimate the area by starting at x n , taking rectangles of width 1 again and
then using the right endpoint as the height of the rectangle. This will give an estimation of the
area under f ( x ) on the interval [ n, ) , as shown in Fig. 2.66(b). Again, we can see that the
remainder, Rn , is again this estimation and in this case it will underestimate the area. Thus
Rn f ( x ) dx . Combining the two inequalities
n
f ( x ) dx Rn f ( x ) dx s n f ( x) dx sn Rn s n f ( x ) dx . Thus
n 1 n n 1 n
sn f ( x) dx s sn f ( x) dx .
n 1 n
Comparison Test: In this case, unlike with the integral test, we may or may not be able to get an
idea of how good a particular partial sum will be as an estimate of the exact value of the series.
Much of this will depend on how the comparison test is used. As in comparison test, given a
series an , we find a second series bn that converged and an bn for all n . Actually we
are estimating the actual value of the series an , whose remainder is Rn i n1 ai , using the
n
partial sum sn i 1 ai . Let the remainder of the series bn is Tn i n 1 bi . As an bn so
Rn Tn . When using the comparison test it is often the case that the bn are fairly nice terms and
that we might actually be able to get an idea on the size of Tn . For instance, if our second series is
a p series we can use the results from above to get an upper bound on Tn as:
Rn Tn g ( x) dx , where g ( n) bn . Also if the second series is a geometric series then we
n
will be able to compute Tn exactly. If we are unable to get an idea of the size of Tn then
comparison test is not helpful.
Alternating Series Test: The integral test and comparison test requires all the terms must be
positive. But in alternating series test there is no restriction for the terms to be positive. Let we
n
have a convergent series an ( 1) bn , which satisfies the condition of alternating series
test, so bn 0 for all n . We want to know how good of an estimation of the actual series value
will the partial sum, sn , be. If the series, an , converges to s will lies between sn and sn 1 for
any n and thus s sn sn 1 sn bn 1 Rn s sn bn 1
Ratio Test: Let the series an is convergent. To get an estimate of the remainder let us first
define the sequence rn an 1 an . We now have two possible cases:
If {rn } is a decreasing sequence and rn 1 1 then Rn an 1 (1 rn 1 )
If {rn } is an increasing sequence then Rn an 1 (1 L) , where L lim an 1 an
n
upon the values of x . Let there is a number R such that the power series will converge for x a R
and will diverge for x a R . This number, R , is called the radius of convergence for the series.
Note that the series may or may not converge if x a R . Also the interval of all x ’s, including the
endpoints if need be, for which the power series converges is called the interval of convergence of
the series. So if we know that the radius of convergence of a power series is R , then we have the
following:
Convergence of power series: a R x a R
Divergence of power series: x a R and x a R
The interval of convergence must then contain the interval a R x a R since we know that the
power series will converge for these values. We also know that the interval of convergence cannot
contain x ’s in the ranges x a R and x a R since we know the power series diverges for these
value of x . Therefore, to completely identify the interval of convergence all that we have to do is
determine if the power series will converge for x a R or x a R . If the power series converges
for one or both of these values then we’ll need to include those in the interval of convergence.
At x a , the power series becomes n0 cn (a a )n n 0 cn (0)n c0 (0)0 n1 cn (0)n c0 and
so the power series converges. Thus it is important to note that no matter what else is happening in the
power series we are guaranteed to get convergence for x a . The series may not converge for any
other value of x , but it will always converge for x a .
Example 2.338 [XE-2008 (2 marks)]: The radius of convergence of the real power series
m0 (m !)2
(2m 1)! x m is
(a) 4 (b) 3 (c) 2 (d) 1
Solution (a): By applying the ratio test, for the convergence of the given power series m1 am ,
( m !) 2 m am 1 (m 1)!2 x m1 (2m 1)!
where am x , we have L lim 1 L lim
(2m 1)! m am m 2(m 1) 1 ! ( m !) 2 x m
( m 1) 2 x (1 1 m) 2 x (1 0) 2 x x x
L lim lim lim
m (2m 3)(2m 2) m (2 3 m)(2 2 m) m (2 0)(2 0) 4 4
Now for the convergence of the given power series we have L 1 (1 4) x 1 x 4 . Thus the
given power series will converge for x 4 and diverge for x 4 and thus the radius of
convergence is R 4 . Also the interval of convergence is x 4 x (4, 4) .
Example 2.339 [CS-1993, EC-1993 (1 mark)]: Find The radius of convergence of the power series
m1(3m)! (m!)3 x3m .
Solution: By applying the ratio test, for the convergence of the given power series m1 am , where
am
(3m)!
x 3 m , we have L lim
am 1 3(m 1) ! x3( m1) (m !)3
3
1 L lim 3
(m !) m am m
(m 1)! (3m)! x3 m
2 m m
Example 2.340 [AE-2014 (1 mark)]: The series s m 1 m 3 ( x 2) converges for all x
with x 2 R given by
(a) R 0 (b) R 3 (c) R (d) R 1 3
Solution (b): By applying the ratio test, for the convergence of the given power series m1 am ,
m2 m am1 (m 1) 2 ( x 2) m 1 3m
where am ( x 2) , we have L lim 1 L lim m 1
3m m am m 3
2
m ( x 2)
m
( m 1) 2 ( x 2) (1 1 m) 2 ( x 2) (1 0) 2 ( x 2) ( x 2)
L lim 2
lim L 1 if
m 3m m 3 3 3
( x 2) 3 1 x 2 3 ; also L 1 if ( x 2) 3 1 x 2 3 . Thus for x 2 3 the given
series converges and for x 2 3 the given series diverges. Thus the radius of convergence R 3 ;
also the interval of convergence is x 2 3 3 x 2 3 1 x 5 .
Power Series and Functions: Every power series can be represented as a function of x . For e.g.,
in the infinite geometric series, n0 ar n a (1 r ) , provided r 1 (the infinite geometric series
converges for r 1 and diverges for r 1 ) if we take a 1 and r x , the series becomes
n0 x n 1 (1 x) , provided x 1 . So we can represent the infinite geometric series as a function
f ( x ) 1 (1 x ) with the power series n 0 x n provided x 1 . We can clearly plug any number
other than x 1 into the function, however, we will only get a convergent power series if x 1 . This
means the equality, in n0 x n 1 (1 x) , will holds if x 1 and for any other value of x the
equality will not hold. Also the radius of convergence of the infinite geometric series is R 1 and the
interval of convergence is x 1 . Thus we can represent many functions as power series and it will be
important to recognize that the representations will often only be valid for a range of x ’s and that
there may be values of x that we can plug into the function that we cannot plug into the power series
representation.
every order. Now we need to determine what are the coefficients cn . Let us first evaluate everything
at x a , so using the power series and its derivatives we find that: f ( a ) c0 , f ( a ) c1 ,
f (a ) 2c2 c2 f (a ) 2 , f ( a ) 3(2)c3 c3 f ( a) {3(2)} , , cn f ( n ) ( a ) n ! . So, a
power series representation for the function f ( x ) about x a is the Taylor series for f ( x ) about
x a is
f ( n ) ( a) n f (a ) 2 f ( a ) 3
f ( x ) n 0 ( x a ) f ( a ) f ( a)( x a) ( x a) ( x a) (2.44)
n! 2! 3!
If we use a 0 , then the Taylor series about x 0 is called a Maclaurin series for f ( x ) , i.e.,
f ( n ) (0) n f (0) 2 f (0) 3
f ( x) n 0 ( x ) f (0) f (0) x (2.45) x x
n! 2! 3!
The Taylor series expansion of functions about x 0 is given in Section 2.3.1 under the heading of
Some important expansion. To determine a condition that must be true in order for a Taylor series to
exist for a function let us first define the nth degree Taylor polynomial of f ( x ) as
n f (i ) (a ) i f ( n ) (a ) n
Tn ( x ) i 0
i!
i0 n ! ( x a) , the nth
( x a ) . Notice that for the full Taylor series
degree Taylor polynomial is just the partial sum for the series. Next the remainder is defined to be
Rn ( x) f ( x ) Tn ( x ) . So the remainder is really just the error between the function f ( x ) and the nth
degree Taylor polynomial for a given n ; with this definition note that we can write the function as
f ( x ) Tn ( x ) Rn ( x ) .
Suppose that f ( x ) Tn ( x ) Rn ( x ) , if lim Rn ( x ) 0 for xa R then
n
f ( n ) (a ) n
f ( x ) n 0 ( x a ) on x a R .
n!
Example 2.341 [CE-2000 (2 marks)]: The Taylor expansion of sin x about x 6 is given by
2 3
1 3 1 3 x3 x5 x7
(a) x x x (b) x
2 2 6 4 6 12 6 3! 5! 7!
3 5 7
1 1 1
(c) x x x x (d) 1 2
6 3! 6 5! 6 7! 6
Solution (a): The Taylor expansion of f ( x ) sin x about x 6 is given as:
2 3
1 1
f ( x) f f x f x f x .
6 6 6 2! 6 6 3! 6 6
As, f ( 6) sin( 6) 1 2 , f ( 6) cos( 6) 3 2 , f ( 6) sin( 6) 1 2 , and so on.
2 3
Thus, f ( x ) (1 2) ( 3 2) x ( 6) (1 4) x ( 6) ( 3 12) x ( 6) .
[Similar question was also asked in EC-2009 (2 marks)]
Example 2.342 [EC-2007 (1 mark)]: For the function e x , the linear approximation around x 2 is
(a) (3 x )e 2 (b) 1 x
(c) 3 2 2 1 2 x e 2 (d) e 2
Solution (a): f ( x) e x f ( x ) e x f ( x ) e x f ( x ) e x , so Taylor series
expansion of f ( x ) about x a is f ( x ) e 2 e 2 ( x 2) (1 2!)e 2 ( x 2) 2 (1 3!)e 2 ( x 2) 3 .
Now neglecting the higher power of ( x 2) , we get
f ( x ) e 2 e 2 ( x 2) f ( x ) (3 x )e 2 .
Example 2.343 [EC-2008 (2 marks)]: In the Taylor series expansion of exp( x ) sin( x) about the
2
point x , the coefficient of x is
(a) exp( ) (b) 0.5 exp( ) (c) exp( ) 1 (d) exp( ) 1
x
Solution (b): The Taylor series expansion of f ( x) e sin x about the point x is given as:
f ( x ) f ( ) f ( )( x ) (1 2!) f ( )( x ) 2 (1 3!) f ( )( x )3 . So the coefficient of
( x ) 2 is f ( ) 2! . As f ( x) e x cos x f ( x ) e x sin x f ( ) e 0 e . Thus the
coefficient of ( x ) 2 is e 2! e 2 0.5e
[Similar question was also asked in ME-2008 (1 mark)]
Example 2.344 [EC-2008 (1 mark)]: Which of the following functions would have only odd powers
of x in its Taylor series expansion about the point x 0 ?
(a) sin( x3 ) (b) sin( x 2 ) (c) cos( x3 ) (d) cos( x 2 )
Solution (a): The Taylor series expansion of f ( x ) sin x and g ( x ) cos x about x 0 is given as:
x3 x5 x2 x4
f ( x ) sin x x and g ( x ) cos x 1 . So
3! 5! 2! 4!
3 3 3 x9 x15 2 2 2 x6 x10
f ( x ) sin( x ) x , f ( x ) sin( x ) x ,
3! 5! 3! 5!
3 3 x6 x12 2 2 x4 x8
g ( x ) cos( x ) 1 , g ( x ) cos( x ) 1
2! 4! 2! 4!
Thus the Taylor series expansion of f ( x 3 ) sin( x 3 ) have only odd powers of x .
k 2 X
Example 2.345 [MT-2011 (2 marks)]: Y k1 1 exp , where k1 , k2 and k3 are
k3 X
constants. If k 2 X k3 X , the value of Y up to first order approximation would be
k 2 X k 2 X k 2 X k 2 X
(a) Y k1 1 (b) Y k1 1 (c) Y k1 (d) Y k1
k3 X k3 X
k3 X k3 X
Solution (d): e x 1 x (1 2!) x 2 (1 3!) x 3 ; now let x ( k 2 X ) ( k3 X ) and for first order
approximation, the higher powers of x should be neglected as k2 X k3 X and so e x 1 x , thus
Y k1 (1 e x ) k1{1 (1 x))} k1 x k1 (k 2 X ) ( k3 X ) .
Example 2.346 [CH-2012 (1 mark)]: For the function f (t ) e t , the Taylor series approximation
for t is
(a) 1 (t ) (b) 1 (t ) (c) 1 (t 2 2 2 ) (d) 1 t
Solution (b): Firstly we have to find the Taylor series expansion of f (t ) e t about t 0 . As
f (t ) e t f (t ) (1 )e t f (t ) (1 2 )e t f (t ) (1 3 )e t and so on. So
2 3
f (0) 1 , f (0) 1 , f (0) 1 , f (0) 1 and so on. Thus
2 3
1 2 1 3 t 1 t 1 t
f (t ) f (0) f (0)(t 0) f (0)(t 0) f (0)(t 0) 1 2
Now
2! 3! 2! 3! 3
for t , the higher powers of t 0 . Thus f (t ) 1 t
Example 2.347 [EC-2014 (1 mark)]: The series n 0 (1 n !) converges to
Example 2.350 [CH-2014 (1 mark)]: If f ( x ) is real and continuous function of x , the Taylor series
expansion of f ( x ) about its minima will NEVER have a term containing
(a) first derivative (b) second derivative (c) third derivative (d) any higher derivative
Solution (a): As for a real valued function y f ( x) , the Taylor series expansion about x a is
f ( x ) f ( a) f ( a )( x a ) (1 2!) f ( a)( x a ) 2 (1 3!) f ( a )( x a)3 . Now for minima at
x a , f ( a) 0 . So f ( x ) f ( a) (1 2!) f ( a)( x a ) 2 (1 3!) f ( a )( x a)3 , which does not
contain the first derivative term.
Example 2.351 [EC-2014 (2 marks)]: The Taylor series expansion of 3sin x 2 cos x is
(a) 2 3 x x 2 (1 2) x3 (b) 2 3 x x 2 (1 2) x3
(c) 2 3 x x 2 (1 2) x 3 (d) 2 3 x x 2 (1 2) x3
Solution (a): The Taylor series expansion of f ( x ) sin x and g ( x ) cos x about x 0 is given as:
sin x x (1 3!) x 3 (1 5!) x 5 ; cos x 1 (1 2!) x 2 (1 4!) x 4 . So the Taylor series expansion
x3 x5 x2 x4 2 x3 x 4
of h( x) 3sin x 2 cos x 3 x 2 1 2 3x x
3! 5! 2! 4! 2 12
b
Improper integral of first kind: A definite integral a f ( x ) dx is called an improper integral of
first kind if the range of integration is not finite (i.e., either a or b or a and b )
1 1 x dx, 0 1 (1 x ) dx,
2 2
and the integrand f ( x) is bounded on [ a, b] . For e.g.,
1 (1 x ) dx, 1 3x
2
(1 2 x )3 dx are improper integrals of first kind. In an improper integral of
first kind, the interval of integration is one of the following types [ a, ) , (, b] , (, ) .
t t b
If a f ( x)dx exists for every t a then, a f ( x) dx lim f ( x ) dx . If t f ( x) dx exists for
t a
b b
every t b then, f ( x) dx lim f ( x ) dx .
t t
or lim
t t
b
f ( x) dx exists finitely and this limit is called the value of the improper integral.
t
If lim f ( x ) dx or lim
t a
t t
b
f ( x) dx is either or , then the integral is said to be divergent.
c c
If f ( x)dx and c f ( x) dx are both convergent then f ( x)dx f ( x)dx c f ( x ) dx ,
where c is any number. The improper integral f ( x)dx is said to be divergent if either of the
right hand side is divergent.
a
Example 2.352 [CE-2000 (1 mark)]: The following integration lim x 4 dx
a 1
a
x 3
( 3)
a
1
lim 1 (3a ) (1 3) 1 3
a
3
Example 2.353 [CH-2012 (2 marks)]: If a is constant, then value of the integral a 2 xe ax dx is
0
2
3. The integral 0 1 (1 cos x) dx, is an improper integral of second kind since integrand
1 (1 cos x) becomes infinite at x [0, 2 ] .
1
4. 0 (sin x) x dx, is a proper integral since lim
x 0
(sin x) x 1 .
If f ( x) is continuous on the interval [a, b) and not continuous at x b then,
b t
a f ( x ) dx lim f ( x ) dx provided the limit exists and is finite. Note that we have to use a LHL
t b a
here as the interval of integration is entirely on the left side of the upper limit.
If f ( x ) is continuous on the interval (a, b] and not continuous at xa then,
b b
a f ( x ) dx lim f ( x ) dx provided the limit exists and is finite. Note that we have to use a RHL
ta t
here as the interval of integration is entirely on the right side of the lower limit.
c b
If f ( x ) is not continuous at x c , where a c b and a f ( x) dx and c f ( x) dx are both
b c b b
convergent then a f ( x) dx f ( x) dx f ( x ) dx . Now the improper integral
a c a f ( x) dx will
converge if both of the improper integral on the right hand side converges. If either of the two
b
integrals on the right hand side is divergent then the improper integral a f ( x) dx diverges.
c b
If f ( x ) is not continuous at x a and x b and if a f ( x) dx and c f ( x) dx are both
b c b
convergent then a f ( x) dx f ( x) dx f ( x ) dx , where a c b . Again this requires both
a c
b
integrals on RHS to be convergent in order for improper integral a f ( x) dx to also be convergent.
Absolute and Conditional convergence: Let the function f ( x) be defined for all x a . If the
integral a f ( x ) dx converges, then the integral a f ( x) dx also converges and is called absolutely
convergent; and in this case a f ( x) dx
a
f ( x) dx . If the integral a f ( x ) dx converges and
a f ( x ) dx diverges, then the integral a f ( x) dx is called conditionally convergent.
Example 2.354 [CS-1993, EC-1993 (1 mark)]: Which of the following improper integrals is/are
convergent?
1 sin x 4 x 1 2(1 cos x )
1 3sin (2 x )
(a) dx (c) dx (d) dx
0 1 cos x
(b) 1 x
dx 0 1 x 2 0 x5 2
1 sin x 1 2 sin( x 2) cos( x 2) 1 cos( x 2)
Solution (d): Option (a) I dx 2
dx dx
0 1 cos x 0 2 sin ( x 2) 0 sin( x 2)
1 1
I 2 ln sin( x 2) 0 2(ln sin(1 2) ln 0) sin x (1 cos x) dx diverges.
0
4 4
Option (b) 0 sin (2 x) 1 1 1 sin (2 x) 2 1 x {1 sin 4 (2 x)} x2 x in
[1, ) . As, 1 1
x dx (2 x )1 , so by comparison test we can say that
1 1 3sin (2 x)
4
x dx diverges as 1
x 1 sin 4 (2 x) x.
x 1 1 x
0 1 x 2 dx 2 ln 1 x 0 ln 1
2 2
Option (c) dx diverges. ln 1
2 1 x2 0
Option (d) As for the given integral we have an improper point x 0 so we have the interval
2 2
1 2(1 cos x ) 1 4 sin ( x 2) 1 sin ( x 2) 1
x (0,1] . As dx 0 x5 2 dx 0 ( x 2)2 x1 2 dx ; and for all x (0,1] , we
0 x5 2
2 sin 2 ( x 2)
have 0 sin ( x 2) 1 ; also lim lim
4 sin 2 ( x 2) x 5 2
1 0 . So using the limit
2
x 0
( x 2) x0 1 x1 2
comparison test we have lim f ( x ) g ( x ) 1 0
x 0
so both
f ( x ) 4 sin 2 ( x 2) x5 2 and
Example 2.355 [CS-2005 (2 marks)]: Let G ( x ) 1 (1 x) 2 i 0 g (i ) x i , where x 1 . What is
g (i ) ?
(a) i (b) i 1 (d) 2 (c) 2i i
2 3 4 5
Solution: 1 (1 x ) 1 x x x x x ; differentiating both sides w.r.t. x , we get
1 (1 x ) 2 1 2 x 3 x 2 4 x 3 5 x 2 i 0 (i 1) x i i 0 g (i ) x i g (i) i 1
1
Option (b)
0 2
x 1
dx tan 1 x 0
tan 1 tan 1 0 2 0 2
Option (c) xe x dx xe x e x ( e e ) ( 0 e 0 e 0 ) 0 (0 1) 1
0 0
1 1 1
Option (d) dx ln 1 x 0 ln(0) ln(1) . So option (d) is unbounded
0 1 x
1 (1 x ) dx is
2
Example 2.357 [ME-2010 (1 mark)]: The value of the integral
(a) (b) 2 (c) 2 (d)
dx b dx b
Solution (d): Let I
1 x 2
lim
a a 1 x 2
lim tan x
a
1
a
lim tan b tan a
a
1 1
b b b
I tan 1 ( ) tan 1 ( ) tan 1 () tan 1 () 2 tan 1 () 2( 2)
[Similar question was also asked in CH-2011 (2 mark)]
1 4 3
Example 2.358 [XE-2010 (1 mark)]: The integral 1 x dx is
(a) an improper integral converging to –6 (b) an improper integral converging to 0
(c) not an improper integral but has value –6 (d) a divergent improper integral
1
Solution (a): Let I x 4 3 dx , then x 0 is the only point of infinite discontinuity of x 4 3 in
1
0 1
0 4 3 1 4 3 3 3 3 3 3 3
[ 1,1] , then I x dx x dx 13
13
13
13
6 . Thus the
1 0 x x 0 ( 1) (1) (0)
1 0
given integral is converging to –6.
Example 2.359: Find at which point the improper integral 0 e x dx is convergent or divergent.
k
Solution: I e x dx lim e x dx I lim[ e x ]k0 lim[ e k e0 ]
0 k 0 k k
k k
I lim(1 e ) 1 0 1 [ lim e e 0]
k k
k
Thus, lim e x dx exists and is finite. Hence the given integral is convergent at 1.
k 0
Exercise: 2.7
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
5. The series n11 3n1 converges to _____.
Which of the following series diverges? (i) n0 (4n n )
2 3
6. (10 2n3 ) ; (ii) n 1 n
(a) Only (i) (b) Only (ii) (c) Both (i) and (ii) (d) Neither (i) not (ii)
7. Which of the following series converges? (i) n0 (4) 3n
) 5 n 1
; (ii) 9 n 1
n 2 n 1
4 . The
series which converges find the value of the series.
(a) (i), 259.2 (b) (ii), 259.2 (c) (i), 50.2 (d) (ii), 50.2
8. If the series n0 1 (n 2
3n 2) converges, then its value is
(a) 1 (b) 0 (c) 2 (d) The series diverges
9. If the series n0 1 (n 2
4n 3) converges, then its value is
(a) 12 5 (b) 5 12 (c) 0 (d) The series diverges
10. Which of the following series converges?
2
n11
(i) n2 1 (n ln n) (ii) n0 ne n (iii) n 4 1 n7 (iv) n
(a) (i) and (ii) (b) (ii) and (iii) (c) (iii) and (iv) (d) (iv) and (i)
11. Which of the following series diverges? (i) n1n 2
( n cos n) ; (ii) 2
n1(n2 2) ( n 2 5)
(a) Only (i) (b) Only (ii) (c) Both (i) and (ii) (d) Neither (i) nor (ii)
n2 (4n2 n)
12. Which of the following series diverges? (i) n0 1 (3 n
n) ; (ii) 3
n 7 n3
(a) Only (i) (b) Only (ii) (c) Both (i) and (ii) (d) Neither (i) nor (ii)
13. Which of the following series converges?
n 2 n 3
( 1) n ( 1) n cos(n )
(i) n 1 2 (ii) n 0 (iii) n2 n
n 5 n 4
(a) (i) and (ii) (b) (iii) and (i) (c) (ii) and (iii) (d) (i), (ii) and (iii)
14. Which of the following series is absolute convergent?
(i) n1(1)n n (ii) n1(1)n2 n2 (iii) n1(sin n) n3
(a) Only (i) (b) Only (ii) (c) Both (1) and (iii) (d) Both (ii) and (iii)
15. Which of the following series diverges?
(i) n1(10)n {42 n 1 ( n 1)} (ii) n0 n! 5n (iii) n 2 n 2 (2n 1)!
(a) Only (i) (b) Only (ii) (c) Only (iii) (d) All (i), (ii) and (iii)
n
16. Which of the following series diverges? (i) n1n n
3 2 n 1
; (ii) (5n 3n )
n 0
3
(7 n 2)
3
(a) Only (i) (b) Only (ii) (c) Both (i) and (ii) (d) Neither (i) nor (ii)
17. The radius of convergence and interval of convergence for the power series given by
n
( 1) n
n1 4n ( x 3)n is, respectively.
(a) 4, 7 x 1 (b) 4, 1 x 7 (c) 3, 7 x 1 (d) 3, 7 x 0
18. The radius of convergence and interval of convergence for the power series given by
n1(2n
n)(4 x 8) n is, respectively.
(a) 1 , 1 x 2 (b) 1 8 , 15 8 x 17 8
(c) 1 4 , 1 x 17 8 (d) 1 8 , 15 8 x 17 8
19. The radius of convergence and interval of convergence for the power series given by
n0 n !(2 x 1)n is, respectively.
(a) 0, x 1 2 (b) 1, x 1 2 (c) 0, x 1 2 (d) 1, x 1 2
20. The radius of convergence and interval of convergence for the power series given by
n1( x 6)n
n n is, respectively.
(a) 0, x (b) 1, x (c) 2, x (d) , x
21. The radius of convergence and interval of convergence for the power series given by
n1 x 2 n (3) n is, respectively.
(a) 3, 3x0 (b) 2, 0 x 3
(c) 3, 3x 3 (d) 2, 3 x0
22. The power series representation of the function f ( x) x (5 x) and its interval of convergence
is, respectively.
(a) 0 xn
5n 1 , x 5 (b) 0 x n1
5n1 , x 5
(c) 0 xn 5n 1 , x 4 (d) 0 x n1 5n1 , x 4
23. The power series representation of the function g ( x ) 1 (1 x ) 2 and its radius of convergence is,
respectively.
(a) n1 nx n 1 , 2 (b) n1 nx n 1 , 1 (c) n1 nx n , 1 (d) n1 nx n , 2
24. The Taylor series expansion for f ( x ) ln x about x 2 is
(a) ln 2 n1 ( 1)
n 1
( n2 n ) ( x 2) n
(b) ln 2 n1 ( 1)
n 1
(2n ) ( x 2) n
(c) ln 2 ( 1)n 1
n 1
n ( x 2) n (d) ln 2 ( 1) n 1
n
( n2 ) ( x 2)
n n
1 1 ( x
1 2 2 1
1 x x 2 1 dx 0 1
2 2
(i) 0 {1 x }dx (ii) ( x 3 8) dx (iii) (iv) 3
x dx
(a) (i) and (iii) (b) (i), (ii) and (iii) (c) (i) and (ii) (d) (iii) and (iv)
Answers Keys
Answer Keys: Exercise: 2.1
1 2 3 4 5 6 7 8 9 10 11 12 13 14 15
b a b 6 a c a c b 6 300 3300 60 45 d
16 17 18 19 20 21 22 23 24 25 26 27 28 29 30
b c a b 32 b b c d a c b 10 160 c
31 32 33 34 35 36 37 38 39 40 41 42 43 44 45
43 c 300 42 c 20 15 b a a d b 6 a c
46 47 48 49 50 51 52 53 54 55 56 57 58
5 d a d 23 a d b d c d d a
[AE – 2017 (1 mark)]: The eigenvalues n and eigenfunctions un ( x ) of the Sturm-Liouville problem
d2y 2
2
k y 0 , 0 x 1 ; y (0) 0 , y (1) 0 are given by:
dx
(a) n n 2 2 ; u n ( x ) sin n x , n 0, 1, 2, ,
(b) n n 2 2 k 2 ; u n ( x) sin kn x, n 0, 1, 2, ,
(c) n n 2 2 k 2 ; u n ( x ) sin n x, n 0, 1, 2, ,
(d) n n 2 2 ; un ( x ) sin n x, n 0, 1, 2, ,
d2y 2 2
Solution (c): Let p 2 , so 2
k p y 0 , whose general solution is given by
dx
y A cos( kpx) B sin( kpx ) . Applying the given conditions: (i) y (0) 0 A 0 ; (ii)
2 2 2 2
y (1) 0 B sin( kp ) 0 kp n p n k ; thus the eigenvalues are n p n k ; and
eigenfunctions are un ( x ) sin( kpx ) sin( n x) , n 0, 1, 2, , .
d2y dy
[AE – 2017 (2 marks)]: Consider the initial value problem: 2
4 6 y f (t ) ; y (0) 2 ,
dt dt
dy st st
1 . If Y ( s ) 0 y (t )e dt and F ( s) 0 f (t )e dt are the Laplace transform of y(t ) and
dt t 0
f (t ) respectively, then Y ( s ) is given by:
F (s) F ( s) 2 s 9 F ( s) F ( s) 2 s 9
(a) 2
(b) 2
(c) 2
(d) 2
( s 4s 6) ( s 4s 6) ( s 4 s 6) ( s 4 s 6)
Solution (b): { y (t )} s 2Y ( s) sy (0) y (0) s 2Y ( s ) 2s 1 ; { y (t )} sY ( s ) y (0) sY ( s ) 2
. So taking Laplace transform on both sides of the given DE, we get:
{ y (t )} 4{ y (t )} 6{ y (t )} { f (t )} s 2Y ( s) 2s 1 4{sY ( s ) 2} 6Y ( s ) F ( s)
F ( s) 2 s 9
Y (s) 2 .
s 4s 6
2 d2y dy
[AE – 2017 (2 marks)]: The equation x 2
5x 4 y 0 has a solution y ( x ) that is:
dx dx
(a) A polynomial in x
(b) Finite series in terms of non-integer fractional power of x
(c) Consists of negative integer powers of x and logarithmic function of x
(d) Consists of exponential functions of x
Solution: The given DE is homogeneous second order Cauchy – Euler equation. Let
y x m dy dx mx m1 d 2 y dx 2 m(m 1) x m 2 ; substituting these in given DE, we get the
auxiliary equation of given DE: m(m 1) 5m 4 0 m 2 4m 4 0 m 2, 2 , i.e. it has
two real and repeated roots. Thus the general solution of given DE is y c1 x 2 c2 x 2 ln x .
[AE – 2017 (2 marks)]: Let u ( x, t ) denote the displacement of a point on a rod. The displacement
2u 2u
satisfies the following equation of motion: 25 0 , 0 x 1 ; with u ( x, 0) 0.01sin(10 x )
t 2 x 2
u
, ( x, 0) 0 ; u (0, t ) 0 , u (1, t ) 0 . The value of u (0.25,1) is _____ (in three decimal places).
t
2u 2 2u
Solution: The solution of the equation c , with u (0, t ) 0 and u ( L, t ) 0 for all t ;
t 2 x 2
u ( x, 0) f ( x ) and (u t ) ( x,0) g ( x) for all 0 xL is given by
L
u ( x, t ) n 1 ( Bn cos n t Bn* sin n t ) sin{( n L) x} , where Bn (2 L ) f ( x) sin{( n L ) x}dx , and
0
L
Bn* {2 (cn )} g ( x ) sin{( n L) x}dx , n cn L , n 1, 2, .
0
u
So we have c 2 25 c 5 ; L 1 ; u ( x, 0) f ( x ) 0.01sin(10 x) ; ( x, 0) g ( x ) 0 . Thus the
t
solutions of the given PDE is for n 10 , we have
* *
u ( x, t ) ( B10 cos 10t B10 sin 10 t ) sin{( L ) x} ( B10 cos 10t B10 sin 10t ) sin( x ) , where
1 1
10 cn L 50 . B10 2 0.01sin(10 x) sin(10 x ) dx 0.01 {1 cos(20 x )}dx 0.01
0 0
* L 1
B10 {2 (cn )} g ( x ) sin{(n L ) x}dx {2 (50 )} 0 sin{10 x}dx 0
0 0
d2y dy
[BT – 2017 (1 mark)]: For y f ( x) , if 2
0, 0 at x 0 , and y 1 at x 1 , the value of
dx dx
y at x 2 is _____.
d2y dy dy
Solution: 2
0 k …(i), where k is a constant; As
0 at x 0 , so from (i), k 0 .
dx dx dx
Now dy dx 0 y c …(ii), where c is a constant. At x 1 , y 1 ; so from (ii), c 1 . Hence
y 1 is the solution of the given DE. Thus at x 2 , y 1 .
[CH – 2017 (2 marks)]: For the initial value problem, dx dt sin t , x(0) 0 , the value of x at
t 3 is _____.
Solution: dx dt sin t dx sin tdt x (t ) cos t c . Applying the given condition, we get
x (0) cos 0 c 0 1 c c 1 . Thus x(t ) cos t 1 ; hence x( 3) cos( 3) 1 0.5 .
[CH – 2017 (2 marks)]: The Laplace transform of a function is [( s 1) {s ( s 2)}] . The initial and
final values, respectively, of the function are
(a) 0 and 1 (b) 1 and 1 2 (c) 1 2 and 1 (d) 1 2 and 0
Solution (b): The initial value of f (t ) is given by:
t 0 s s
lim f (t ) lim{sF ( s)} lim s
s 1
lim
s 1
lim
1 (1 s)
s( s 2) s s 2 s 1 (2 s )
1
t s0 s 0
s 1
The final value of f (t ) is given by: lim f (t ) lim{sF ( s )} lim s lim
s 1
s ( s 2) s 0 s 2 2
1
.
[EC – 2017 (1 mark)]: A periodic signal x (t ) has a trigonometric Fourier series expansion
x (t ) a0 n 1 (an cos n0t bn sin n0 t ) . If x(t ) x ( t ) x{t ( 0 )} , we can conclude that
(a) an are zero for all n and bn are zero for n even
(b) an are zero for all n and bn are zero for n odd
(c) an are zero for n even and bn are zero for n odd
(d) an are zero for n odd and bn are zero for n even
Solution (a): As x(t ) x ( t ) x (t ) is having odd symmetry; and for odd symmetry we have:
2
a0 0 , an 0 and bn
T 0
T
f (t ) sin( n t )dt .
As x(t ) x{t ( 0 )} x{t ( (2 2T ))} x{t T } x (t ) is having half wave symmetry;
and for half wave symmetry:
0, n is even 0, n is even
a0 0 ; an T ; bn T .
(2 T ) f (t ) cos( nt ) dt , n is odd
0 (2 T ) f (t ) sin( nt ) dt , n is odd
0
Hence option (a) is correct.
[EC – 2017 (2 marks)]: Which one of the following is the general solution of the first order
differential equation dy dx ( x y 1) 2 , where x , y are real? (where c is a constant)
[EC – 2017 (2 marks)]: Let x (t ) be a continuous time periodic signal with fundamental period T 1
second. Let {ak } be the complex Fourier series coefficients of x (t ) , where k is integer valued.
Consider the following statements about x (3t ) :
I. The complex Fourier series coefficients of x (3t ) are {ak } where k is integer valued
II. The complex Fourier series coefficients of x (3t ) are {3ak } where k is integer valued
III. The fundamental angular frequency of x (3t ) is 6 rad/s
For the three statements above, which one of the following is correct?
(a) only II and III are true (b) only I and III are true (c) only III is true (d) only I is true
Solution (b): Fourier series coefficient ak is unaffected by scaling operation. Thus statements (I) is
true and statement (II) is false. The time period for x (t ) is T 1 second and if it is compressed by 3,
then the resultant period T 1 3 second. Hence fundamental frequency 2 T 2 (1 3) 6 rad/s.
So statement (III) is correct.
d2y dy
[EC – 2017 (1 mark)]: The general solution of the differential equation 2
2 5 y 0 in terms
dx dx
of arbitrary constants K1 and K 2 is
( 1 6 ) x
(a) K1e K 2 e ( 1 6)x
(b) K1e
( 1 8 ) x
K 2 e ( 1 8)x
( 2 6 ) x ( 2 6 ) x ( 2 8 ) x ( 2 8 ) x
(c) K1e K2e (d) K1e K 2e
Solution (a): We have homogeneous linear equation with constant coefficients, so its solution will be
of the form y e mx y me mx y m 2 e mx , substituting all these in given DE, we get:
e mx ( m 2 2m 5) 0 , as e mx 0 , so m 2 2m 5 0 m 1 6 , which are real and distinct. So
( 1 6 ) x
the general solution of the given DE is K1e K 2 e ( 1 6)x
, where K1 and K 2 are constants.
[EE – 2017 (2 marks)]: Consider the differential equation (t 2 81)( dy dt ) (5t ) y sin(t ) with
y (1) 2 . There exists a unique solution for this differential equation when t belongs to the interval
(a) ( 2, 2) (b) ( 10,10) (c) ( 10, 2) (d) (0,10)
dy 5t sin t
Solution: 2
y 2
…(i), which is a first order linear equation; so
dt t 81 t 81
5t 5
t 2 81dt ln( t 2 81)
I.F. e e (t 2 81) 5 2 .
2
dy 2 5t sin t 2
Thus multiplying (i) with I.F. we get: (t 81)5 2 2 (t 2 81)5 2 y 2 (t 81)5 2
dt t 81 t 81
d
dt
y (t 2 81)5 2 sin t (t 2 81)3 2 y (t 2 81)5 2 {(sin t )(t 2 81) 3 2 }dt c
2
{(sin t )(t 81)3 2 }dt c
y 2 52
.
(t 81) (t 81)5 2
2
If t 9 , then the solution does not exist; and t 9 or t 9 or both lie in the options given in (b),
(c) and (d). So we can say that option (a) is correct.
[EE – 2017 (2 marks)]: Let a causal LTI system be characterized by the following differential
d2y dy dx (t )
equation, with initial rest condition 2
7 10 y (t ) 4 x(t ) 5 where, x (t ) and y (t ) are
dt dt dt
the input and output respectively. The impulse response of the system is { u (t ) is the unit step
function}
(a) 2e 2t u (t ) 7e 5t u (t ) (b) 2e 2t u (t ) 7e 5t u (t )
(c) 7e 2t u (t ) 2e 5t u (t ) (d) 7e 2 t u (t ) 2e 5t u (t )
Solution (b): Taking Laplace transform on both sides of given causal LTI system, we get
d 2y dy dx (t )
2
7 10 y (t ) 4 x(t ) 5
dt dt dt
s 2Y (s ) sY (0) y(0) 7 sY ( s ) y (0) 10Y ( s ) 4 X ( s ) 5 sX ( s ) x(0)
As the initial rest condition Y (0) 0 , y (0) 0 , y (0) 0 , x(0) 0 . So
Y ( s) 4 5s 4 5s 2 7
s 2Y ( s ) 7 sY ( s ) 10Y ( s ) 4 X ( s ) 5sX ( s ) 2
X (s) s 7 s 10 ( s 2)( s 5) s2 s 5
Y (s) 2 7
H (s) . Taking inverse Laplace transform will give h(t ) (impulse response).
X ( s) s2 s5
Thus h(t ) 1
s2
2
s5
1
7
2e 2 t u (t ) 7e 5t u (t ) .
[MA – 2017 (1 mark)]: If x (t ) and y (t ) are the solutions of the system dx dt y and dy dt x
with the initial conditions x (0) 1 and y (0) 1 , then x( 2) y ( 2) equals _____.
dx d 2x dy d 2x
Solution: y x x 0 , whose solution is x (t ) A sin t B cos t .
dt dt 2 dt dt 2
dx(t )
Applying the given conditions, we get x(0) 1 1 B ; y (t ) y (0) 1
dt
A cos 0 B sin 0 1 A 1 .
Thus x(t ) sin t cos t x( 2) 1 . Similarly, y(t ) sin t cos t y ( 2) 1 .
So x( 2) y( 2) 1 1 0 .
d2y
[ME – 2017 (1 mark)]: The differential equation 16 y 0 for (a) no solution
dx 2 (b) exactly two solutions
y( x) with two boundary conditions ( dy dx ) x 0 1 and (c) exactly one solution
( dy dx ) x 1 has (d) infinitely many solutions
2
Solution (a): Given DE is homogeneous with constant coefficient, so substituting
y e mx y me mx y m 2 e mx in given DE, we get: e mx ( m 2 16) 0 ; as e mx 0 , so
( m 2 16) 0 m 0 4i . So the general solution of the given DE is
(0) x
ye (c1 cos 4 x c2 sin 4 x ) (c1 cos 4 x c2 sin 4 x ) y 4( c1 sin 4 x c2 cos 4 x )
Thus ( dy dx ) x 0 1 4c2 1 c2 1 4 ; and ( dy dx ) x 2
1 4c2 1 c2 1 4 . As
value of c2 cannot be both 1 4 and 1 4 . So with given boundary conditions, the given DE has no
solution.
y e (0) x (c1 cos 3 x c2 sin 3 x) (c1 cos 3 x c2 sin 3 x ) y 3( c1 sin 3x c2 cos 3 x )
Thus y (0) 0 c1 0 ; and ( dy dx) x 0 2000 3c2 2000 c2 2000 3 . So at x 1,
y (2000 3) sin 3 94.08 .
[MT – 2017 (2 marks)]: For the second order linear ordinary differential equation,
d2y dy
2
p qy 0 , the following function is a solution: y e x . Which one of the following
dx dx
statements is NOT TRUE?
(a) has two values: one complex and one real (b) 2 p q 0
(c) has two real values (d) has two complex values
Solution (a): Substituting y e y e y 2e x in given DE, we get:
x x
d3y d2y dy
[PE – 2017 (1 mark)]: The roots of the equation 3
6 2
11
6 y 0 are:
dx dx dx
(a) 1, 1, 2 (b) 1, 2, 3 (c) 1, 3, 4 (d) 1, 2, 4
Solution: Given DE is homogeneous with constant coefficient, so substituting
y e mx y me mx y m 2 e mx y m3e mx in given DE, we get:
e mx ( m3 6m 2 11m 6) 0 ; as e mx 0 , so ( m3 6m 2 11m 6) 0 m 1, 2,3 , which are the
roots of auxiliary equation.
[PH – 2017 (2 marks)]: Consider the differential equation ( dy dx) y tan x cos x . If y (0) 0 ,
y ( 3) is _____. (up to two decimal places).
[PE – 2017 (2 marks)]: The temperature time profile for a system is given as follows:
dT
5T 500 , where T is temperature in oC, and t is time in hours. The initial conditions are
dt
T (0) 500 oC. The temperature of the system after 1 hour is _____ oC. (write answer with two
decimal places).
dT d
After multiplying IF to the given DE, we get: e5t 5T (e5t ) 500(e5t ) {Te5 t } 500e5 t
dt dt
5t 5t 5t 5t 5 t
d {Te } 500 e dt Te 100e c T (t ) 100 ce .
Applying the given condition T (0) 500 oC, we get: 500 100 ce 0 c 400
Thus T (t ) 100 400e 5t . So T (1) 100 400e 51 102.69
t s 1 s 1
Thus {e cos 2t} F ( s 1) 2 2
2
.
( s 1) 2 s 2s 5
dy
[TF – 2017 (2 marks)]: If a solution curve of the differential equation x y 2 x 3 passes through
dx
the point (1, 0) , then this curve also passes through the point
(a) ( 1, 0) (b) (0, 1) (c) (2,10) (d) ( 2, 6)
dy dy 1 1
y 2 x 2 …(i), whose I.F. e
( 1 x ) dx 1
Solution (a): x y 2 x3 e ln x eln x . On
dx dx x x
multiplying the I.F.
with (i), we get:
1 dy 1 d y y 2
2 y 2x 2 x d ( y x ) 2 xdx x c …(ii).
x dx x dx x x
0
Thus at (1, 0) from (ii), we get 12 c c 1 .
1
y
So the curve is given by: x 2 1 y x 3 x , which also passes through ( 1, 0) .
x
[XE – 2017 (2 marks)]: If u ( x, t ) g (t ) sin x is the solution of the wave equation: utt u xx , t 0 ,
0 x ; with the initial conditions: u ( x, 0) 2 sin x , ut ( x, 0) 0 , 0 x ; and the boundary
conditions u (0, t ) u ( , t ) 0 , t 0 , then the value of g ( 3) is _____.
Solution: Comparing the given equation utt u xx , with Eq. 3.109, we have c 1 , and given
conditions with Eqs. 3.110 and 3.111, we have L , f ( x ) 2 sin x and g ( x ) 0 . Thus the
solution of the given PDE is from Eq. 3.121, i.e. for n 1 , we have
* *
u ( x, t ) ( B cos t B sin t ) sin{( L ) x} ( B cos t B sin t ) sin x , where c L 1 1
and from Eqs. 3.123 and 3.125,
B
2
L
L
0 f ( x ) sin{( L ) x}dx
2
0 2 sin
2
xdx
2
0 (1 cos 2 x)dx
2
x
sin 2 x
2 0
2
L
B* {2 (cL )} g ( x ) sin{( L ) x}dx (2 ) 0 sin xdx 0 .
0 0
[XE – 2017 (1 mark)]: Consider the ordinary differential equation y y y 0 , where and
are constants. If y ( x ) xe x is a solution of the above equation, then the value of is _____.
Solution: y ( x ) xe x y ( x ) e x xe x y ( x ) 2e x xe x . On substituting these in the given DE,
we get: {2e x xe x } {e x xe x } {xe x } 0 e x {(2 ) x (1 )} 0 ; as e x 0 , so
{(2 ) x(1 )} 0 , which is satisfied if 2 and 1 0 1 . So 3 .
An equation containing the derivatives of one or more dependent variable w.r.t. one or more
independent variables is said to be a differential equation.
If a differential equation contains only ordinary derivatives of one or more dependent variable
w.r.t. a single independent variable, it is said to be an ordinary differential equation (ODE). For
e.g. ( dy dx ) 5 y e x , ( d 2 y dx 2 ) ( dy dx ) 6 y 0 , ( dx dt ) ( dy dt ) 2 x 2 y are ODEs.
An equation involving partial derivatives of one or more dependent variables of two or more
independent variables is called a partial differential equation (PDE). For e.g.
( 2u x 2 ) ( 2u y 2 ) 0 , ( 2 u x 2 ) ( 2 u t 2 ) 2(u t ) , u y v x are PDEs.
Order and Degree of a Differential Equation: The order of a differential equation (either ODE
or PDE) is the order of the highest differential coefficient occurring in it. The degree of a differential
equation which is expressed or can be expressed as a polynomial in the derivatives is the degree of the
highest order derivative occurring in it after it has been expressed in a form free radicals and fractions
as far as derivatives are concerned. The order and degree of a differential equation is a positive
integer. Thus the differential equation (3.1) is of m th order and p th degree.
m m p m 1 m 1 q
f ( x, y ) d y dx ( x, y ) d y dx 0 (3.1)
Example 3.4 [CE-2010 (1 mark)]: The order and degree of the differential equation
( d 3 y dx3 ) 4 ( dy dx )3 y 2 0 are, respectively
(a) 3 and 2 (b) 2 and 3 (c) 3 and 3 (d) 3 and 1
3
Solution (a): The given differential equation can be written as (d y dx ) 16 (dy dx) y
3 2
3 2
; the
highest differential coefficient (hdc) is d 3 y dx 3 so its order is 3; also degree of hdc is 2.
d 3 y dx3
Example 3.5: Find the order and degree of the differential equation e x (d 2 y dx 2 ) y 0 .
d 3 y dx3 2 2
Solution: e x( d y dx ) y 0 Clearly, the differential equation is of 3rd order, but its
degree is not defined as it cannot be written as a polynomial equation in derivatives, and hence it
cannot be expressed as polynomial of derivatives.
Example 3.6 [TF-2011 (1 mark)]: The order and degree of the following differential equation are
x(dy dx) 2 ( dy dx) y 2
(a) order 1, degree 1 (b) order 1, degree 2 (c) order 2, degree 1 (d) order 2, degree 2
Solution (b): The given differential equation can be written as x( dy dx ) y 2 ( dy dx ) 2 0 and
2
Linear and Non-Linear Differential Equation: A linear differential equation is any differential
equation that can be written in the following form:
an ( x ) y ( n ) ( x ) an 1 ( x) y ( n 1) ( x ) a1 ( x ) y( x ) a0 ( x ) y ( x ) g ( x ) (3.2)
In linear differential equation there are no products of the function, y ( x ) , and its derivatives and
neither the function or its derivatives occur to any power other than the first power. The coeffiecients
a0 ( x), a1 ( x ), an ( x ) and g ( x ) can be zero or non-zero functions, constants or non-constant
functions, linear or non-linear functions. So only the function, y ( x ) , and its derivatives are used in
determining if a differential equation is linear. For e.g., ( y x )dx 4 xy dy 0 , y 2 y y 0 ,
x 3 ( d 3 y dx 3 ) x ( dy dx ) 5 y e x , etc., are linear first, second and third order ODE.
In Eq. 3.2, the dependent variable y and all derivatives y , y , , y ( n ) are of the first degree.
The coefficients a0 , a1 , , an of y , y , , y ( n ) depends at most on the independent variable x .
Two important cases of Eq. 3.2 are given as:
Linear first – order ( n 1 in Eq. 3.2) a1 ( x)( dy dx ) a0 ( x ) y g ( x ) (3.3)
Linear second – order ( n 2 in Eq. 3.2) a2 ( x)( d 2 y dx 2 ) a1 ( x )( dy dx) a0 ( x ) y g ( x) (3.4)
A nonlinear ordinary differential equation is simply one that is not linear. Nonlinear functions of the
dependent variable or its derivatives, such as sin y or, cannot appear in a linear equation. For e.g., the
non - linear term: non - linear term: non - linear term:
coefficient depends on y
nonlinear function of y
power not 1
2 2
equations, (1 y ) y 2 y e x , ( d y dx ) sin y 0 , ( d 4 y dx 4 ) y2 0,
etc., are non – linear first, second and fourth order ODEs.
Example 3.11 [CH-2014 (2 marks)]: (a) non-linear differential equation of first degree
The differential equation (b) linear differential equation of first degree
( d y dx ) x ( dy dx ) x y e x is a
2 2 2 3 (c) linear differential equation of second degree
(d) non-linear differential equation of second degree
Solution (b): The given differential equation is of the form given by Eq. 3.4 so it is a linear
differential equation. Also, the higest differential coefficient is d 2 y dx 2 , which is of second order,
whose degree in the given equation is 1. So the given differential equation is a linear differential
equation of second order and first degree.
Example 3.13 [MN-2012 (1 mark)]: The 2nd order differential equation having a solution
y ( A x) B , where A and B are constants, is
Formation of Partial Differential Equation: Partial differential equation can be formed either
by elimination of arbitrary constants or by the elimination of arbitrary functions from a relation
involving three or more variables. We will discuss the above mentioned concept with the help of
following examples:
Example 3.14 [CE-2010 (2 marks)]: The partial differential equation that can be formed from
z ax by ab has the form with p z x and q z y
(a) z px qy (b) z px pq (c) z px qy pq (d) z qy pq
Solution (c): Differentiating both sides of the given equation partially w.r.t. x and y , we get
z x a p (given) and z y b q (given), respectively. Putting these result in the given
equation we get z px qy pq .
represented by the general solution when taken together will give the locus of the differential
equation. Since there is one arbitrary constant in the general solution of the equation of first order, the
locus of the equation can be said to be made up of single infinity of curves.
A first order and first degree differential equation can be written as f ( x, y ) dx g ( x, y )dy 0 or
dy dx f ( x, y ) g ( x, y) or dy dx ( x, y ) , where f ( x, y ) and g ( x, y ) are obviously the functions
of x and y . It is not always possible to solve this type of equations. The solution of this type of
differential equations is possible only when it falls under the category of some standard forms.
Example 3.16 [CE-2006 (2 marks)]: A spherical naphthalene ball exposed to the atmosphere loses
volume at a rate proportional to its instantaneous surface area due to evaporation. If the initial
diameter of the ball is 2 cm and the diameter reduces to 1 cm after 3 months, the ball completely
evaporates in
(a) 6 months (b) 9 months (c) 12 months (d) Infinite time
3
Solution (a): As the rate of volume, V (4 3) r , of spherical ball is proportional to its
instantaneous surface area, 2
S 4 r , so
dV dt kS (d dt ) (4 3) r 3 k (4 r 2 )
4 r 2 ( dr dt ) k (4 r 2 ) dr dt k dr kdt . Integrating both sides we get
Example 3.17 [CE-2007 (2 marks)]: A body originally at 60oC cools down to 40oC in 15 min. when
kept in air at a temperature of 25oC. What will be the temperature of the body at the end of 30 min.?
(a) 35.2oC (b) 31.5oC (c) 28.7oC (d) 15oC
Solution (b): As from Newton’s law of cooling the rate of temperature ( T ) is directly proportional to
the temperature difference ( T T0 ) dT dt k (T T0 ) 1 (T T0 ) dT kdt . Integrating both
kt
sides we get 1 (T T0 ) dT k dt ln(T T0 ) kt ln c T T0 ce . Now from the data
k 0
given we have, T0 25 oC and at t 0 , T 60 oC 60 25 ce 35 ; also at t 15 minutes,
o k 15 15 k
T 40 C 40 25 35e e 3 7. Now at t 30 minutes, we have
30 k 15 k 2 2 o o
T 25 35e 25 35(e ) 25 35(3 7) 31.428 C 31.5 C.
Example 3.18 [CH-2007 (1 mark)]: The initial condition for which the following equation
( x 2 2 x )( dy dx ) 2( x 1) y ; y ( x0 ) y0 has infinitely many solutions, is
(a) y ( x 0) 5 (b) y ( x 0) 1 (c) y ( x 2) 1 (d) y ( x 2) 0
Solution (d): ( x 2 2 x )( dy dx ) 2( x 1) y (1 y ) dy 2( x 1) ( x 2 2 x ) dx . Integrating both
sides we get, (1 y )dy 2( x 1)
( x 2 2 x) dx ln y ln( x 2 2 x) ln c y c ( x 2 2 x ) . Now
applying the conditions given in option (a) and (b) gives 5 0 and 1 0 , respectively, which is not
possible so option (a) and (b) are not correct as we get no solution for the given condition. For the
condition given in option (c), we get c 1 8 , thus y ( x 2 2 x ) 8 which is an unique solution. For
the condition given in option (d) we get 0 c 0 c can be any number thus we get only infinitely
many solutions; thus option (d) is correct.
Example 3.20 [ME-2007 (2 marks)]: The solution of dy dx y 2 with initial value y (0) 1 is
bounded in the interval
(a) x (b) x 1 (c) x 1, x 1 (d) 2 x 2
Solution (c): dy dx y 2 (1 y 2 )dy dx . Integrating both sides
2
(1 y ) dy dx (1 y ) x c . Now applying the given condition we get, 1 1 0 c
c 1 1 y x 1 y 1 ( x 1) . Hence the solution is valid for x 1 x 1, x 1 .
Example 3.21 [IN-2008 (2 marks)]: Consider the differential equation dy dx 1 y 2 . Which one of
the following can be a particular solution of this differential equation?
(a) y tan( x 3) (b) y tan x 3 (c) x tan( y 3) (d) x tan y 3
Solution (a): dy dx 1 y 2 1 (1 y 2 ) dy dx . Integrating both sides 1 (1 y 2 ) dy dx
tan 1 y x c y tan( x c ) . Now from the given options we can say that option (a) is correct.
[Similar questions were also asked in CE-1999, CE-2004, ME-2003, PI-2010, ME-2011 (2
marks), AE-2012, AG-2008 (1 mark)]
Example 3.23 [MT-2009 (2 marks)]: The solution function y f ( x) for the ordinary differential
equation dy dx 3 x 2 2 x , passes through (1,1) . The magnitude of y at x 3 is
(a) 0 (b) 18 (c) 19 (d) 21
Solution (c): dy dx 3 x 2 2 x dy (3 x 2 2 x )dx . Integrating both sides dy (3 x 2 2 x ) dx
y x3 x 2 c . Applying the given condition, we get 1 1 1 c c 1 . Thus
3 2 3 2
y x x 1 y x 3 3 3 1 19 .
Example 3.24 [AE-2011 (2 marks)]: The solution of dy dt y 3et t 2 with initial condition y (0) 1
is given by
1
(a) et (t 3) 2 9 4et 1
(b) (c) 2
(d)
9 t 2
5 2e (t 2t 2) (t 2) t 2
5 2e (t 2t 2)
Solution (b): dy dt y 3et t 2 (1 y 3 ) dy et t 2 dt . Integrating both sides (1 y 3 ) dy et t 2 dt . As
3 2
LHS of the integral (1 y )dy y 2 ; and for RHS of the integral, let ‘ t 2 ’ as the 1st, ‘ e t ’ as
the 2nd function et t 2 dt t 2 et 2tet dx t 2et 2 tet 1 et dt t 2et 2 te t et et (t 2 2t 2) .
Thus y 2 2 e t (t 2 2t 2) (5 2) 1 y 2 5 2et (t 2 2t 2) y 1 5 2e (t 2t 2)
t 2
Example 3.25 [CH-2011 (2 marks)]: Which ONE of the following choices is a solution of the
differential equation dy dx ( y 2 x ) ( y x ) (2 x ) ? Note c is a real constant.
c x2 c 2 x2 c x3 c 2 x3
(a) y (b) y (c) y (d) y
c 2 x2 c x2 c 2 x3 c x3
Solution (d): The given differential equation can be written in variable separable form as:
dy dx dy dx 1 1 1 dx 1 1 1 dx
2
dy dy .
y y2 x ( y 2)( y 1) x 3 y 1 y 2 x 3 y 2 y 1 x
1 y2 3 y 2 3 1 2 k 3 x3
ln( y 2) ln( y 1) ln x ln c ln ln( kx ) ( kx ) y 3 3
.
3 y 1 y 1 1 k x
Now let k 3 1 c so y (c 2 x 3 ) (c x 3 ) .
Example 3.26 [CE-2012 (2 marks)]: The solution of the ordinary differential equation
( dy dx) 2 y 0 for the boundary condition, y 5 at x 1 is
(a) y e 2x (b) y 2e 2 x (c) y 10.95e 2 x (d) y 36.95e 2 x
Solution (d): ( dy dx) 2 y 0 (1 y ) dy 2dx . Now integrating both sides, (1 y ) dy 2 dx
ln y 2 x ln c ln( y c) 2 x y c e 2 x y ce 2 x . Applying boundary conditions, we
get 5 ce 2 c 5e 2 36.95 , thus the solution of the give ODE, with given boundary condition, is
y 36.95e 2 x
[Similar questions were also asked in CE-2007, CE-2008, EC-2011, EE-2005, ME-1994 (1
mark), CE-2004, TF-2012, CH-2013 (2 marks)]
Example 3.27 [EE-2014 (1 mark)]: A particle, starting from origin at t 0 second, is travelling
along x axis with velocity v ( 2) cos{( 2)t} m/s. At t 3 second, the difference between the
distance covered by the particle and the magnitude of displacement from the origin is ………
Solution: v ( 2) cos{( 2)t} dx dt ( 2) cos{( 2)t} dx ( 2) cos{( 2)t}dt .
Integrating both sides we get dx ( 2) cos{( 2)t}dt x sin{( 2)t} c . As at t 0 , we
have x 0 0 sin ( 2) 0 c c 0 . So at t 3 seconds, the position of the particle is
x sin(3 2) . From the graph of x sin t between t 0 and t 3 2 , the distance of the particle
from the origin is 3m and displacement of the paricle is –1m, which is of magnitude 1m. So the
difference between the distance covered by the particle and the magnitude of displacement from the
origin is 3 1 2 .
Example 3.28 [ME-2014 (1 mark)]: The solution of the initial value problem dy dx 2 xy ;
y (0) 2 is
2 2 2 2
(a) 1 e x (b) 2e x (c) 1 e x (d) 2e x
Solution (b): dy dx 2 xy (1 y )dy 2 xdx . Now integrating both sides (1 y ) dy 2 xdx
2
ln y x 2 ln c y ce x . So from the given options we can say that option (b) is in the form
of our result.
[Similar questions were also asked in EC-2008, EE-2011 (1 mark), MT-2010 (2 marks)]
Example 3.30 [AG-2016 (1 mark)]: The general solution of the differential equation
dy
e3 x 2 y x 2 e 2 y is ___.
dx
(a) C (1 2)e 2 y (1 3)(e 3 x x 3 ) (b) C e 2 y (1 3)(e3 x x 2 )
(c) C (1 3) e 2 y (1 2)(e 3 x x 2 ) (d) C e 2 y (1 3)(e3 x x 3 )
dy
Solution (a): e3 x 2 y x 2e2 y e2 y (e3 x x2 ) e 2 y dy (e3 x x 2 )dx …(i). On integrating both
dx
2y 3x 2 e2 y e3 x x3
sides of (i), we get e dy (e x ) dx C C (1 2)e 2 y (1 3)(e3 x x 3 ) .
2 3 3
Example 3.31 [ME-2014 (2 marks)]: The general solution of the differential equation
dy dx cos( x y ) , with c as a constant, is
(a) y sin( x y ) x c (b) tan ( x y ) 2 y c
(c) cos ( x y ) 2 x c (d) tan ( x y ) 2 x c
Solution (d): Let x y t 1 ( dy dx ) dt dx dy dx ( dt dx ) 1 . So dy dx cos( x y )
dt dx 1 cos t dt dx 1 cos t 2 cos 2 (t 2) sec 2 (t 2) dt 2dx . Now integrating both
2
sides we get, sec (t 2) dt 2 dx 2 tan(t 2) 2 x c tan ( x y ) 2 x c , where c c 2 .
Solution of Homogenoeus Differential Equation: A function f ( x, y ) is called a homogeneous
function of degree n if f ( x, y ) n f ( x, y ) . For example, f ( x, y ) x 2 y 2 3xy is a
homogeneous function of degree 2, because f ( x , y ) 2 x 2 2 y 2 3( x )( y ) 2 f ( x, y ) . A
homogeneous function f ( x, y ) of degree n can always be written as f ( x, y ) x n f ( y x) or
f ( x , y ) y n f ( x y ) ; thus a differential equation of the form dy dx f ( x, y ) is homogeneous if the
function f ( x, y ) depends only on the ratio x y or y x . [This point was asked in ME-1995 (1
mark)]. If a first-order first degree differential equation is expressible in the form
dy dx f ( x, y ) g ( x, y) where f ( x, y ) and g ( x, y ) are homogeneous functions of the same degree,
then it is called a homogeneous differential equation. Such type of equations can be reduced to
variable separable form by the substitution y v x . The algorithm for solving homogeneous
differential equation is given as:
1. Put the differential equation in the form dy dx ( x, y) ( x, y)
2. Put y v x dy dx v x ( dv dx ) in the equation in step 1 and cancel out x from the right
hand side. The equation reduces to the form v x (dv dx ) F (v ) .
3. Shift v on RHS and separate the variables v and x
4. Integrate both sides to obtain the solution in terms of v and x
5. Replace v by y x in the solution obtained in step 4 to obtain the solution in terms of x and y .
Equation reducible to homogeneous form: A first order, first degree non – homogeneous
differential equation is of the form
dy dx ( a1 x b1 y c1 ) ( a2 x b2 y c2 ) , a1 a2 b1 b2 (3.10)
It can be reduced to homogeneous form by certain substitutions. Put x X h, y Y k , where h
and k are constants, which are to be determined. Since, dy dx ( dy dY )(dY dX )( dX dx ) dY dX .
Substituting this value in Eq. 3.12, we have,
dY dX ( a1 X b1Y ) a1h b1k c1 ( a2 X b2Y ) a2 h b2 k c2 (3.11)
Now h, k will be chosen such that a1h b1k c1 0 and a2 h b2 k c2 0 , i.e.,
h (b1c2 b2 c1 ) k (c1a2 c2 a1 ) 1 ( a1b2 a2b1 ) . For these values of h and k the Eq. 3.13 reduces to
dY dX ( a1 X b1Y ) (a2 X b2Y ) which is a homogeneous differential equation and can be solved
by the substitution Y v X . Replacing X and Y in the solution so obtained by x h and y k
respectively, we can obtain the required solution in terms of x and y .
Solution of Differential Equation of the form of Leibnitz’s Linear Equation: The general
form of a linear differential equation of first order (known as Leibnitz’s linear equation) is given as,
( dy dx) P y Q (3.12)
where P and Q are functions of x (or constants). [This point was asked in CE-1997 (1 mark)]
For example, ( dy dx ) xy x3 , x (dy dx) 2 y x 3 , ( dy dx ) 2 y sin x etc. are linear differential
equations. The following are the steps for solving Leibnitz’s Linear Equation:
1. Write the differential equation in the form ( dy dx ) Py Q and obtain P and Q .
4. Integrate both sides of the equation obtained in step 3 w.r.t. x to obtain y (I.F.) Q (I.F.) dx c . This
gives the required solution.
Example 3.33 [CH-2007 (2 marks)]: The solution of the following differential equation
x (dy dx ) y ( x 2 1) 2 x 3 is
(c) c1 x c2 x 2
2 2
(a) 0 (b) 2 ce x 2 (d) 2 x cxe x 2
Solution (d): The given differential equation can be written as ( dy dx ) x (1 x ) y 2 x 2 …(i),
which is in Leibnitz’s linear form, in which P( x ) x (1 x ) and Q ( x) 2 x 2 . So
x (1 x ) dx
I.F. e
2 2 2
2 ln x )
e( x (e x 2 ) (eln x ) (e x 2 ) x ; multiplying both sides of (i) by the I.F.,
2 2 2 2
ex
2
1 dy x2 2ex
d ex 22
x2 2
ex 2 2
x y 2 xe y 2 xe d y 2 xe x 2 dx . Now
x dx x x dx x x
2 2 2
integrating both sides (e x 2
x ) y 2e x 2
c y 2 x cxe x 2
.
I.F. e e
Pdx 3 (1 x ) dx 3
e 3ln x e ln x 1 x 3 . Now multiplying both sides of (i) by the I.F., we get
1 dy 3 1 d 1 1 1 1 y 1
3
4
y 2
3 y 2 d 3 y 2 dx 3 c . Now applying the
x dx x x dx x x x x x x
Q ( x ) cosec x . So I.F. e
cot xdx
e ln sin x sin x ; multiplying both sides of the given differential
equation by the I.F., sin x ( dy dx ) y cos x 1 ( d dx)( y sin x ) 1 d ( y sin x ) dx ; integrating
both sides we get, d ( y sin x) dx y sin x x c y ( x c) cos ecx .
Example 3.38 [CH-2014 (2 marks)]: The integrating factor for the differential equation
( dy dx ) { y (1 x )} (1 x ) is
(a) 1 (1 x ) (b) 1 x (c) x(1 x) (d) x (1 x )
Solution (a): Comparing the given differential equation with Eq. 3.12, we have
P ( x ) 1 (1 x) , so the I.F. e
{1 (1 x )}dx ln(1 x ) ln{1 (1 x )}
e e 1 (1 x)
[Similar questions were also asked in AG-2012 (1 mark), TF-2014 (2 marks)]
Example 3.39 [CE-2014 (1 mark)]: The integrating factor for the differential equation
( dP dt ) k 2 P k1 L0 e k1t is
(a) e k1t (b) e k2 t (c) e k1t (d) e k2 t
Solution (d): Comparing the given differential equation with Eq. 3.12, I.F. e
k 2 dt
e k 2t
Example 3.40 [EC-2014 (1 mark)]: Which one of the following is a linear non-homogeneous
differential equation, where x and y are the independent and dependent variables respectively?
(a) ( dy dx) xy e x (b) ( dy dx ) xy 0 (c) ( dy dx ) xy e y (d) ( dy dx ) e y 0
Solution (a): Differential equation given in option (a) is a first order linear non-homogeneous
equation; but differential equation given in option (b) is a first order linear homogeneous equation.
The differential equations given in option (c) and (d) are first order non-linear equations.
Example 3.41 [XE-2016 (1 mark)]: Let y ( x ) be the solution of the initial value problem
dy
2 xy x ; y (0) 0 . Find the value of lim y ( x ) .
dx x
g ( y ) N ( x, y ) ( y ) M ( x, y ) dx (3.19)
It is important to realize that Eq. 3.19 is independent of x because
N N M
N ( x, y ) M ( x, y ) dx M ( x, y ) dx 0.
x y x y x x y
4. Now integrate Eq. 3.19 with respect to y and substitute the result in Eq. 3.18. The implicit
solution of the equation is f ( x, y ) c .
d (e y x ) ( xe y dy e y dx ) x 2 d (e x y ) ( ye x dx e x dy ) y 2
d ( x 2 y 2 ) ( xdx y dy ) x2 y 2 d ( x m y n ) x m1 y n 1 ( mydx nxdy )
{1 (1 n)}d [ f ( x, y )]1 n f ( x, y ) ( f ( x, y )) n d{1 ( xy )} ( xdy ydx ) ( x 2 y 2 )
d (1 2) log{( x y ) ( x y )} ( x dy y dx ) ( x 2 y 2 )
Example 3.44 [CH-2008 (1 mark)]: Which ONE of the following is NOT an integrating factor for
the differential equation xdy ydx 0 ?
(a) 1 x 2 (b) 1 y 2 (c) 1 ( xy) (d) 1 ( x y )
Solution (d): For option (a): ( xdy ydx ) x 2 0 d ( y x ) 0 which can be integrated easily
Example 3.45 [EC-2012, EE-2012, IN-2012 (1 mark)]: With initial condition x(1) 0.5 , the
solution of the differential equation, t ( dx dt ) x t is
(a) x t (1 2) (b) x t 2 (1 2) (c) x t 2 2 (d) x t 2
Solution: The given differential equation can be written as ( d dt )(t x ) t d (t x ) t dt . Integrating
both sides d (tx) t dt tx (1 2)t 2 c . Now applying the given condition, we get,
1 0.5 (12 2) c c 0 . So the solution is t x t 2 2 x t 2 .
[Similar question was also asked in ME-2009 (2 marks)]
Solution of Non – exact Differential equation: In previous section, for solving linear
differential equation we multiply the equation by integrating factor. The same basic idea works for a
non-exact differential equation M ( x, y ) dx N ( x, y ) dy 0 ; i.e., it is possible to find an integrating
factor ( x, y ) so that after multiplying it by the non-exact differential equation, the resulting equation
becomes an exact differential equation, i.e.,
( x, y ) M ( x, y ) dx ( x, y ) N ( x, y ) dy 0 (3.20)
is an exact differential equation. In an attempt to find ( x, y ) , we turn to the criterion (Eq. 3.18) for
exactness. Eq. 3.20 is exact if and only if ( M ) y ( N ) x , where the subscripts denote partial
derivatives. By the product rule of differentiation the criterion for exactness becomes,
M y y M N x x N x N y M (M y N x ) (3.21)
For determining the only unknown ( x, y ) , we must solve a partial differential equation. Since we
are not prepared to do that, we make a simplifying assumption. Suppose is a function of one
variable; for example, say that depends only on x . In this case, x d dx and y 0 , so that
Eq. 3.21 can be written as,
d dx ( M y N x ) N (3.22)
We are still at an impasse if the quotient ( M y N x ) N depends on both x and y . However, if after
all obvious algebraic simplifications are made, the quotient ( M y N x ) N turns out to depend solely
on the variable x , then Eq. 3.22 is a first-order ordinary differential equation. We can finally
( M y N x ) N dx
determine because Eq. 3.22 is separable as well as linear, as ( x) e . In like manner,
it follows from Eq. 3.21 that if depends only on the variable y , then
d dy ( N x M y ) M (3.23)
In this case, if ( N x M y ) M is a function of y only, then we can solve Eq. 3.23 for .
We now summarize the results for the non-exact differential equation
M ( x, y ) dx N ( x, y ) dy 0 (3.24)
If ( M y N x ) N is a function of x alone, then an integrating factor for Eq. 3.24 is
( M y N x ) N dx
( x) e (3.25)
If ( N x M y ) M is a function of y alone, then an integrating factor for Eq. 3.24 is
( N x M y ) M dy
( y) e (3.26)
( y) e
3 (1 y ) dy 3
e 3ln y e ln y y 3 .
x ’ alone, so I.F. e e
{( M y N x ) N }dx {1 x}dx
e ln x x .
d 2 y dx 2 f ( x ) (3.27)
n
Eq 3.27 may be re-written as ( d dx)( dy dx) f ( x) d ( dy dx) f ( x) dx . On integrating we get,
( dy dx ) f ( x ) dx c1 ( dy dx ) F ( x) c1 , F ( x ) f ( x ) dx c1dx (3.28)
Eq. 3.28 can be written as, dy f ( x) dx c1dx and on integrating it we get,
y F ( x) dx c1 x c2 H ( x ) c1 x c2 , H ( x) F ( x )dx (3.29)
where, c1 , c2 are arbitrary constants.
Example 3.50 [CE-2001 (2 marks)]: The solution of the differential equation d 2 y dx 2 3 x 2 with
boundary conditions y (0) 2 and y (1) 3 is
(a) y (1 3) x 3 (1 2) x 2 3x 6 (b) y 3 x 3 (1 2) x 2 5 x 2
(c) y (1 2) x 3 x 2 (1 2)5 x 2 (d) y x 3 (1 2) x 2 5 x (3 2)
Solution (c): d 2 y dx 2 3 x 2 ( d dx)( dy dx ) 3x 2 d ( dy dx) (3x 2)dx . Integrating both
sides we get d (dy dx ) (3 x 2)dx dy dx (3 2) x 2 2 x c1 . Now applying the given
condition y (1) 3 to our result, we get 3 (3 2)12 2 1 c1 c1 5 2
dy dx (3 2) x 2 2 x (5 2) dy (3 2) x 2 2 x (5 2) dx . Again integrating both sides we
dy (3 2) x
2 3 2
get, 2 x (5 2) dx y (1 2) x x (5 2) x c2 . Again applying the given
condition y (0) 2 , we get 2 (0 3 2) 0 2 (5 2) 0 c2 c2 2 . Thus our final solution is
y (1 2) x3 x 2 (5 2) x 2
[Similar question was also asked in PI-2009 (2 marks)]
(a) x (1 3) x 3
2
(b) x (1 3!) x 3 (c) x (1 3!) x 3 (d) x 2 (1 3!) x 3
Solution (c): (d dx ) b( x )
a(x)
f ( x, t ) dt f x, b( x ) b( x ) f x, a( x ) a ( x )
b( x)
a( x)
f x ( x, t ) dt
which is the Leibniz integral rule. So applying it in the given equation we get
x
( d dx ) ( x) 1 sin( x x) ( x ) ( d dx) x sin( x 0) (0) (d dx )(0) cos( x ) ( ) d
0
x
( x ) 1 cos( x ) ( ) d . Again applying th Leibniz integral rule to our result we get
0
2 2 x
(d dx ) ( x ) 0 cos( x x ) ( x ) ( d dx )( x ) cos( x 0) (0) ( d dx )(0) sin( x ) ( ) d
0
x
( x ) ( x ) 0 sin( x ) ( ) d x d ( d dx ) ( x ) xdx integrating it we get
0
Exercise: 3.1
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
x x x x 1
(a) 2 x 2c (b) 2 x 2c (c) x 2c (d) cx 2
log z log z log z log z 2
14. The slope of the tangent at ( x, y ) to a curve passing through 1, 4 as given by
2
( y x ) cos ( y x ) , then the equation of the curve is
(a) y tan 1 log(e x ) (b) y x tan 1 log( x e)
(c) y x tan 1 log(e x ) (d) None of these
2 2 2
15. The solution of the equation x ( d y dx ) ln x when x 1 , y 0 and dy dx 1 is
1 1 1 1
(a) (ln x) 2 ln x (b) (ln x) 2 ln x (c) (ln x ) 2 ln x (d) (ln x ) 2 ln x
2 2 2 2
2
16. A continuously differentiable function ( x ) in (0, ) satisfying y 1 y , y (0) 0 y ( ) is
(a) tan x (b) x( x ) (c) ( x )(1 e x ) (d) None of these
17. The order of the differential equation whose general solution is given by
x c5
y (c1 c2 ) cos( x c3 ) c4e where c1 , c2 , c3 , c4 , c5 are arbitrary constants, is _____.
18. If y (t ) is a solution of (1 t )(dy dt ) ty 1 and y (0) 1 , then y (1) is equal to
(a) 1 2 (b) e 1 2 (c) e 1 2 (d) 1 2
Boundary value problem (BVP): Another type of problem consists of solving a linear differential
equation of order two or greater in which the dependent variable y or its derivatives are specified at
different points. A problem such as,
Solve: a2 ( x)( d 2 y dx 2 ) a1 ( x )( dy dx) a0 ( x ) y g ( x) (3.32)
Subject to: y (a ) y0 , y (b) y1 (3.33)
is called a boundary value problem (BVP). The prescribed values y ( a ) y0 and y (b) y1 are called
boundary conditions. A solution of these problems is a function satisfying the differential equation on
some interval I , containing a and b , whose graph passes through the two points (a, y0 ) and (b, y1 ) .
For a 2nd order differential equation other pairs of boundary conditions could be y (a ) y0 , y (b) y1 ;
y ( a) y0 , y(b) y1 ; y (a ) y0 , y (b) y1 , where y 0 and y1 denote arbitrary constants.
A BVP can have Many, One, or No solutions even if it satisfy all the condition of Existence
of a unique solution. For e.g., the two-parameter family of solutions of the differential equation
x 16 x 0 is x c1 cos 4t c2 sin 4t .
(a) Suppose we wish to determine the solution of the equation that further satisfies the boundary
conditions x(0) 0, x( 2) 0 . x (0) 0 c1 0 and x( 2) 0 c2 . Hence the
boundary-value problem x 16 x 0, x (0) 0, x( 2) 0 has infinitely many solutions.
(b) If the BVP in (a) is changed to x 16 x 0, x (0) 0, x( 8) 0 then x (0) 0 c1 0 ,
but x ( 8) 0 c2 0 . Hence the BVP has a unique solution, i.e., x 0 .
(c) If we change the BVP to x 16 x 0, x (0) 0, x( 2) 1 then x (0) 0 c1 0 but
x ( 2) 1 x c2 sin 2 1 c2 0 c2 , so the BVP has no solution.
Differential Operators: In calculus differentiation is often denoted by the capital letter D , i.e.,
dy dx Dy . The symbol D is called a differential operator because it transforms a differentiable
function ( y ) into another function. Higher-order derivatives can be expressed in terms of D in a
natural manner: ( d dx )( dy dx) (d 2 y dx 2 ) D( Dy ) D 2 y and in general, d n y dx n D n y .
Solution of Non – Homogeneous Equations: Let y p ( x ) be any particular solution of the non –
homogeneous linear nth order differential Eq. 3.35 on an interval I , and let y1 , y2 , , yn be a
fundamental set of solutions of the associated homogeneous differential Eq. 3.34 on I . Then the
general solution of the equation on the interval is y c1 y1 ( x ) c2 y2 ( x ) cn yn ( x ) y p ( x ) , where
c1 , c2 , , cn are n arbitrary constants.
Thus the general solution of a non – homogeneous linear equation consists of the sum of two
functions y c1 y1 ( x ) c2 y2 ( x ) cn yn ( x ) y p ( x ) yc ( x ) y p ( x ) .
The linear combination yc ( x) c1 y1 ( x ) c2 y2 ( x ) cn yn ( x ) , which is the general solution of
Eq. 3.34 is called the complementary function for Eq. 3.35
In other words, to solve a nonhomogeneous linear differential equation, we first solve the
associated homogeneous equation and then find any particular solution of the nonhomogeneous
equation. The general solution of the nonhomogeneous equation is then
y complementary function any particular solution yc y p .
Superposition Principle for Non – Homogeneous Equations: Let y p1 , y p2 , , y pk be k
particular solutions of the non – homogeneous linear nth order differential Eq. 3.35 on an interval
I corresponding, in turn, to k distinct functions g1 , g 2 , , g k . That is, suppose y pi denotes a
particular solution of the corresponding differential equation
an ( x )( d n y dx n ) an 1 ( x )( d n 1 y dx n 1 ) a1 ( x )( dy dx ) a0 ( x ) y g i ( x ) (3.36)
where i 1, 2, , k . Then,
y p y p1 ( x ) y p2 ( x ) y pk ( x ) (3.37)
is a particular solution of
dny d n 1 y dy
an ( x ) n an 1 ( x ) n 1 a1 ( x) a0 ( x ) y g1 ( x ) g 2 ( x ) g k ( x ) (3.38)
dx dx dx
y2 y1 ( x ) 1 y12 ( x ) e
P ( x ) dx
dx (3.42)
Note that, Eq. 3.42 which is a solution of Eq. 3.40, y1 and y 2 are linearly independent on any interval
on which y1 is not zero.
Example 3.53 [TF-2009 (2 marks)]: The general solution of the differential equation
x 2 ( d 2 y dx 2 ) 2 x ( dy dx ) ( y x 2 ) 0 is
x
(a) y c1 cos x c2 sin x (b) y c1e x c2 e x (c) y c1 (1 x )c2 e1 x (d) y (c1 c2 x)e
Solution (c): As the third term on the LHS of the given DE is y x 2 , so the general solutions given in
options (a), (b) and (c) are not correct. Now we will check for the general solution given in option (c).
If y c1 (1 x )c2 e1 x is the general solution of the given differential equation then y1 e1 x is its
2 1x
solution. So y1 (1 x )e y1 e1 x {(1 2 x ) x 4 } ; substituting these in the given DE we get
LHS x2 e1 x {(1 2 x) x4 } 2 x( e1 x x 2 ) (e1 x x 2 ) (e1 x x2 )(1 2 x 2 x 1) 0 RHS y1 e1 x
is a solution of the given differential equation. Hence we can say that option (c) is correct.
Example 3.54 [XE-2009 (2 marks)]: Let y1 ( x ) and y2 ( x ) be two linearly independent solutions of
( d 2 y dx 2 ) (6 x)( dy dx ) q ( x) y 0 , x (1,3) , where q( x ) is a continuous function in (1,3) . If the
Wronskian of y1 ( x ) and y2 ( x ) at x 1 , denoted by w( y1 , y2 )(1) , is 1, then w( y1 , y2 )(2) is
(a) 1 26 (b) 1 23 (c) 1 2 (d) 2
Solution (a): If y1 ( x ) and y2 ( x ) are two linearly independent solutions of y p ( x ) y q ( x ) y 0 ,
where p ( x) 1 x ; then we have y1 p ( x ) y1 q( x ) y1 0 …(i) and y2 p ( x) y2 q( x ) y2 0 …(ii).
Now multiplying (i) by y 2 and (ii) by y1 , we get y2 y1 p ( x) y2 y1 q( x ) y2 y1 0 …(iii) and
y1 y2 p ( x) y1 y2 q( x ) y1 y2 0 …(iv). So (iv) – (iii) y1 y2 y2 y1 p( x )( y1 y2 y2 y1 ) 0 . As
W ( y1 , y2 ) y1 y2 y2 y1 W ( y1 , y2 ) y1 y2 y2 y1 . So we have W p ( x )W 0 W W p ( x) ;
now integrating both sides w.r.t. x, we get
6 6
(W W )dx p( x)dx ln W (6 x)dx 6 ln x ln c ln(c x ) W c x . Now as
W y1 (1), y2 (1) 1 c 16 1 c 1 W y1 ( x), y2 ( x) 1 x 6 W y1 (2), y2 (2) 1 26 .
(1 x 2 )}dx ln(1 x 2 ) e
P ( x ) dx 2
As P ( x)dx {2 x e ln(1 x ) 1 (1 x 2 ) . Given y1 x is the
solution of given equation, and let the second solution be Q ( x ) y2 so from Eq. 3.42,
y2 x (1 x 2 )e
P ( x ) dx
dx . y2 x (1 x 2 )e
P ( x ) dx
dx x 1 x 2 (1 x 2 ) dx .
1 1 1 1 1 1 1 1 1 1 x
2 2
dx 2
dx ln 1 x ln 1 x ln .
x (1 x ) x 2(1 x ) 2(1 x ) x 2 2 x 2 1 x
y2 1 ( x 2) ln (1 x ) (1 x ) and so y ( x ) c1 x c2 ( x 2) ln 1 x 1 x 1 is the
general solution.
Auxiliary Equation: We begin by considering the special case of the second order equation
ay by cy 0 (3.44)
where a , b and c are constants. To find a solution of the form y emx , then after substitution of
y me mx and y m 2 e mx , Eq. 3.44 becomes e mx (am 2 bm c ) 0 . e mx 0 for all x , hence
2
am bm c 0 (3.45)
Eq. 3.45 is called the auxiliary equation of the differential equation 3.44. Since, the two roots of Eq.
2
3.45 are m1 , m2 (b b 4ac ) 2a there will be three forms of the general solution of Eq. 3.44
corresponding to the three cases:
Case I: Distinct Real Roots: If b 2 4ac 0 then Eq. 3.45 has two real and distinct roots m1 and
mx m x
m2 , and we have two solutions y1 e 1 and y2 e 2 . We see that these functions are linearly
independent on (, ) and hence form a fundamental set. It follows that the general solution of
Eq. 3.44 on this interval is
mx m x
y c1e 1 c2 e 2 (3.46)
Case II: Repeated Real Roots: If b 2 4ac 0 then Eq. 3.45 has two real and equal roots, i.e.,
m1 m2 m b 2a , we necessarily obtain only one exponential solution y1 e mx . From Eq.
3.42, the second solution of the equation 3.44 is given as,
ay by cy 0 y (b a ) y (c a ) y 0 P b a 2 m and Qc a
y2 e mx e 2 mx e 2 mx dx e mx dx xe mx . Hence the general solution of Eq. 3.44 is
m1 x m1x m1x
y2 c1e c2 xe (c1 c2 x )e (3.47)
2
Case III: Conjugate Complex Roots: If b 4ac 0 then m1 and m2 are complex, hence we
can write m1 , m2 i , where i 2 1 , , 0 and are real. Formally, there is no difference
between this case and Case I, and hence y c1e ( i ) x c2 e ( i ) x is the general solution.
However, in practice we prefer to work with real functions instead of complex exponentials by
using Euler’s formula: ei cos i sin , where is any real number. It follows from this
formula that ei x cos x i sin x and e i x cos x i sin x ei x e i x 2 cos x and
ei x e i x 2i sin x . Since y c1e ( i ) x c2 e ( i ) x is a solution of Eq. 3.44 for any choice of
constants c1 and c2 , the choices c1 c2 1 and c1 1, c2 1 gives the two solutions:
Example 3.59 [ME-2005 (2 marks)]: Which of the following is a solution of the differential equation
( d 2 y dx 2 ) p (dy dx ) ( q 1) y 0 ?
Example 3.61 [IN-2007 (2 marks)]: The boundary value problem y y 0 , y (0) y ( ) 0 will
have non-zero solutions if and only if the values of are
(a) 0, 1, 2, (b) 1, 2, 3, (c) 1, 4, 9, (d) 1, 9, 25,
Solution (c): As we have homogeneous linear equation with constant coefficient so putting
y e mx y me mx y m 2 e mx in the given DE we get e mx ( m 2 ) 0 . As e mx 0 , so we
have the auxiliary equation m 2 0 whose roots are m 0 i which is a conjugate complex
root case, so 0 and . So the general solution of given DE is
0
y e (c1 cos x c2 sin x ) c1 cos x c 2 sin x . Now we have to find the constants c1 and
c2 . As y (0) 0 c1 0 and y ( ) 0 c2 sin 0 sin 0 as we have to find non-zero
solutions of y , so c2 0 . Thus sin 0 m m ., where m I , i.e.,
m {0, 1, 2, } . So m 2 {1, 4, 9, } .
Example 3.65 [PI-2010 (1 mark)]: Which one of the following differential equations has a solution
given by the function y 5sin 3x ( 3) ?
dy 5 dy 5 d2y d2y
(a) cos(3 x ) 0 (b) cos(3 x) 0 (c) 9 y 0 (d) 9y 0
dx 3 dx 3 dx 2 dx 2
Solution (c): Twice differentiating the given solution we get
2 2
dy d y y d y
15 cos 3x 2 45 sin 3 x 45 9 y 2 9 y 0
dx 3 dx 3 5 dx
Example 3.67 [IN-2013 (2 marks)]: The maximum value of the solution y (t ) of the differential
equation y (t )
y (t ) 0 with initial conditions y (0) 1 and y (0) 1 , for t 0 is
(a) 1 (b) 2 (c) (d) 2
Solution: As we have homogeneous linear equation with constant coefficient so putting
y e mx y me mx y m 2 e mx in the given DE we get e mx ( m 2 1) 0 so we have the auxiliary
equation m 2 1 0 as e mt 0 . As the discriminant of the auxiliary equation
D (0) 2 4 1 1 4 0 so we have conjugate complex roots case; also the roots of our auxiliary
equation are m1 , m2 0 i 0 and 1 , so the general solution of given DE is
y (c1 cos x c2 sin x ) y ( c1 sin x c2 cos x) . Now applying the given condition y (0) 1
1 (c1 c2 0) c1 1 ; also at y (0) 1 1 ( c1 0 c2 ) c2 1 . So the general solution is
y (cos x sin x ) whose maximum value is 12 12 2.
Example 3.68 [ME-2013, PI-2013 (2 marks)]: The solution of the differential equation
( d 2 u dx 2 ) k ( du dx ) 0 , where, k is a constant, subjected to the boundary conditions u (0) 0 and
u ( L) U is
x 1 e kx 1 e kx 1 e kx
(a) u U (b) u U kL
(c) u U kL
(d) u U kL
L 1 e 1 e 1 e
Solution (b): As we have homogeneous linear equation with constant coefficient so putting
mx mx 2 mx
u e u me u m e in the given DE we get e mx ( m 2 km) 0 so we have the
auxiliary equation m 2 km 0 (as e mt 0 ), m1 0 and m2 k so we have distinct real root
case; so the general solution of given DE is u c1e 0 c2 e kx c1 c2 e kx . Now applying the given
condition u (0) 0 0 c1 c2 e 0 c1 c2 0 ; also u ( L ) U U c1 c2 e kL ; solving these two
equations we get, c2 U (e kL 1) and c1 U (e kL 1) . So the general solution is
y U (e kL 1) Ue kx (e kL 1) U (1 e kx ) (1 e kL ) .
[Similar question was also asked in MT-2016 (1 mark)]
Example 3.69 [EC-2014 (1 mark)]: Solution (a): As we have homogeneous linear equation with
If the characteristic equation of the constant coefficient so putting y e mx y me mx
differential equation
2 2 y m 2 e mx in the given DE we get e mx (m 2 2 m 1) 0
( d y dx ) 2 ( dy dx ) y 0 has
two equal roots, then the values of so we have the auxiliary equation m 2 2 m 1 0 as
mx
are e 0 . It is given that the auxiliary equation has two
(a) 1 (b) 0, 0 equation roots so the discriminant of auxiliary equation must
(c) j (d) 1 2 be zero, i.e., 4 2 4 0 2 1 1 .
Example 3.70 [EE-2016 (1 mark)]: A function y (t ) , such that y (0) 1 and (a) 5e 1
d2y dy (b) 5e 2
y (1) 3e 1 , is a solution of the differential equation 2 y 0 .
dt 2 dt (c) 7e 1
Then y (2) is (d) 7e 2
Solution (b): We have homogeneous linear equation with constant coefficients, so its solution will be
of the form y e mx y me mx y m 2 e mx , putting all these in given DE, we get
mx
e mx ( m 2 2m 1) 0 (m 2 2m 1) 0 , as e 0 ; so ( m 2 2m 1) 0 ( m 1) 2 0 m 1
Thus we have real and equal roots and so the general solution of given DE is
y (c1 c2 x )e mx (c1 c2 x)e x .
Example 3.71 [BT-2016 (2 marks)]: ( d 2 y dx 2 ) y 0 . The initial conditions for this second order
homogeneous DE are y (0) 1 and dy dx 3 at x 0 . The value of y when x 2 is _____.
Solution: We have homogeneous linear equation with constant coefficients, so its solution will be of
the form y e mx y me mx y m 2 e mx , putting all these in given DE, we get
mt
e mx ( m 2 1) 0 ( m 2 1) 0 , as e 0 ; so ( m 2 1) 0 m 1,1 . Thus we have real and
distinct roots and so the general solution of given DE is y c1e x c2 e x y c1e x c2 e x . Now
y (0) 1 1 c1 c2 …(i); and y (0) 3 3 c1 c2 …(ii). From (i) and (ii), we have c1 1
and c2 2 . Hence y ( x ) e x 2e x ; so y (2) e 2 2e 2 14.642 .
Higher Order Differential Equations: In general, to solve an nth -order differential equation
3.43, we must solve an nth -degree polynomial equation (obtained by putting y emx in Eq. 3.43)
an m n an 1m n 1 a2 m 2 a1m a0 0 (3.49)
If all the roots of Eq.3.39 are real and distinct, then the general solution of Eq. 3.43 is
m1x m2 x mn x
y c1e c2 e cn e (3.50)
It is somewhat harder to summarize the analogues of Cases II and III because the roots of an
auxiliary equation of degree greater than two can occur in many combinations. For example, a
fifth-degree equation could have five distinct real roots, or three distinct real and two complex
roots, or one real and four complex roots, or five real but equal roots, or five real roots but two of
them equal, and so on. When m1 is a root of multiplicity k of an nth -degree auxiliary equation
(that is, k roots are equal to m1 ), it can be shown that the linearly independent solutions are
m1 x m1 x m1x m1 x
e , xe , x 2e , , x k 1e and the general solution must contain the linear combination
m1 x mx mx mx
c1e c2 xe 1 c3 x 2 e 1 c k x k 1e 1 (3.51)
Finally, it should be remembered that when the coefficients are real, complex roots of an auxiliary
equation always appear in conjugate pairs. Thus, for example, a cubic polynomial equation can
have at most two complex roots.
Example 3.76 [ME-1996 (2 marks)]: The particular solution for the differential equation
( d 2 y dx 2 ) 3( dy dx ) 2 y 5 cos x is
(a) 0.5cos x 1.5sin x (b) 1.5cos x 0.5sin x (c) 1.5sin x (d) 0.5cos x
Solution (a): As g ( x ) 5 cos x so the particular solution is in the form of y p A cos x B sin x . As
y A sin x B cos x and y A cos x B sin x ; substituting these in the given DE we get,
( A cos x B sin x ) 3( A sin x B cos x) 2( A cos x B sin x) 5 cos x . Comparing the coefficients
of sin x and cos x of LHS with sin x and cos x of RHS, we get A 3B 2 A 5 A 3B 5
and B 3 A 2B 0 B 3 A 0 ; solving these two we get A 0.5 and B 1.5 . So the particular
solution is y p 0.5 cos x 1.5sin x .
[Similar questions were also asked in ME-2006, XE-2008, TF-2008, TF-2013 (2 marks)]
get, 16 A cos 2 x 16B sin 2 x A cos 2 x B sin 2 x 15cos 2 x ; now comparing the coefficients of
cos 2x and sin 2x , we get 15 A 15 A 1 and 15B 0 B 0 y p cos 2 x and thus the
x
general solution is y yc y p c1e c2e x c3 cos x c4 sin x cos x .
[Similar question was also asked in CE-2000 (5 marks)]
Example 3.78 [ME-2001 (5 marks)]: Solve the differential equation ( d 2 y dx 2 ) y x with the
following conditions: (1) at x 0, y 1 ; (2) at x 2 , y 2 .
Solution: We first solve the associated homogeneous equation y y 0 . From the quadratic
formula we find that the roots of the auxiliary equation m2 1 0 are m1 , m2 0 i . Hence the
complementary function is yc e 0 x (c1 cos x c2 sin x) c1 cos x c2 sin x . As we have linear
function of g ( x ) x , so the particular solution is of the form y p Ax B y p A y p 0 so
substituting these in the given DE we get, 0 Ax B x A 1 and B 0 ; so y p x and thus
the general solution is y yc y p c1 cos x c2 sin x x . Now applying the given condition
y (0) 1 1 c1 ; y ( 2) 2 2 c2 2 c2 0 . So the general solution is y cos x x
Example 3.79 [EE-2005 (2 marks)]: For the equation x(t ) 3x (t ) 2 x (t ) 5 , the solution x (t )
approaches the following values at t
(a) 0 (b) 5 2 (c) 5 (d) 10
Solution (b): The associated homogeneous equation for the given DE is x 3 x 2 x 0 . From the
quadratic formula for the auxiliary equation m 2 3m 2 0 , its roots are m1 1, m2 2 . Hence
the complementary function is xc c1e t c2 e 2 t . Now, because the function g (t ) is constant, let us
assume a particular solution that is also constant as, x p A . We seek to determine the coefficient A
for which x p is a solution of the given equation. Substituting x p and the derivatives xp 0 and
xp 0 into the given DE, we get, x 3 x 2 x 5 0 0 2 A 5 A 2.5 . Thus a particular
Example 3.83 [XE-2010 (2 marks)]: Which one of the following is a particular solution of the
ordinary differential equation x( d 2 y dx 2 ) (dy dx) 2 x 2 f ( x ) ?
(a) x 2 x f ( x) dx x 3 f ( x ) dx (b) x 2 f ( x ) dx x 2 f ( x ) dx
(c) x 2 x f ( x ) dx x 3 f ( x ) dx (d) x 2 f ( x ) dx x 2 f ( x ) dx
Solution (d): The given differential equation can be written as
2
d y dy 2 d dy dx dy dx
x 2 x 2 f ( x) 2 f ( x) d 2 f ( x) dx
dx dx dx x x
dy dx 2 x f ( x )dx dy 2 x f ( x) dx dx dy 2 x f ( x )dx dx
y2 f ( x)dx xdx 2 (d dx ) f ( x)dx xdx dx x f ( x)dx x
2 2
f ( x ) dx
Example 3.86 [AE-2014 (2 marks)]: Solution to the boundary value problem 9( d 2 u dx 2 ) u 5 x ,
0 x 3 with u (0) 0 , (du dx) x 3 0 is
(a) u ( x) {15e (1 e 2 )}(e x 3 e x 3 ) 5 x (b) u ( x) {15e (1 e 2 )}(e x 3 e x 3 ) 5 x
(c) u ( x ) 15sin( x 3) cos(1) 5 x (d) u ( x) 15 sin( x 3) cos(1) (5 54) x3
Solution (a): As 9( d 2 u dx 2 ) u 5 x (i) non-homogeneous DE whose associated homogeneous
DE is 9(d 2u dx 2 ) u 0 …(ii), which has constant coefficients. So putting u e mx u me mx
u m e
2 mx
in (i), we get e mx (9m 2 1) 0 . As e mx 0 , so we have the auxiliary equation
2
9m 1 0 whose roots are m 1 3 which is a distinct real root case. So the general solution of
(ii) is u c c1e x 3 c2 e x 3 . Now, as the function g ( x ) 5 x so the choice of the particular solution is
u p Ax b u p A u p 0 . Now putting these into the given DE we get 0 ( Ax B) 5 x
Example 3.87 [MA-2014 (2 marks)]: The solution to the initial value problem
( d 2 y dt 2 ) 2( dy dt ) 5 y 3e t sin t , y (0) 0 and ( dy dt ) x 0 3 is
(a) y (t ) et (sin t sin 2t ) (b) y (t ) e t (sin t sin 2t ) (c) y (t ) 3et sin t (d) y (t ) 3e t sin t
Solution (b): The given DE is ( d 2 y dt 2 ) 2( dy dt ) 5 y 3e t sin t …(i) whose associated
2 2
homogeneous is ( d y dt ) 2( dy dt ) 5 y 0 …(ii) which is solved by putting
mx mx 2 mx mx 2 mx
ye y me y m e , so e (m 2m 5) 0 as e 0 so the auxiliary equation is
2
m 2m 5 0 m 1 2i we have conjugate complex root case. So the complementary
solution of (i) is given as: yc e t (c1 cos 2t c2 sin 2t ) . Now as g (t ) e t sin t so the choice of our
particular solution is y p e t ( A cos t B sin t ) yp e t {( A B ) cos t ( A B ) sin t}
y p e t {2 B cos t 2 A sin t} . So putting these in (i) we get,
e t (2 B cos t 2 A sin t ) 2{( A B) cos t ( A B ) sin t} 5( A cos t B sin t ) 3e t sin t
t
e t {(3 A) cos t (3B ) sin t} 3e t sin t 3B 3 B 1 and 3 A 0 A 0 . So y p e sin t .
Hence y yc y p e t (c1 cos 2t c2 sin 2t ) e t sin t . Now applying the given conditions:
y (0) 0 0 e 0 (c1 0) 0 c1 0 y e t (c2 sin 2t sin t ) ; now as
t 0
y e {c2 sin 2t sin t 2 c2 cos 2t cos t} y (0) 3 3 e (2 c2 1) c2 1 . Hence the
t
general solution of the given differential equation is y e (sin 2t sin t ) .
Example 3.88 [CE-2016 (2 marks)]: The respective expression for complimentary function and
d4y d2y 2
particular integral part of the solution of the differential equation 4
3 2
108 x are
dx dx
(a) c1 c2 x c3 sin 3 x c4 cos 3x and 3x 4 12 x 2 c
(b) c2 x c3 sin 3 x c4 cos 3 x and 5 x 4 12 x 2 c
(c) c1 c3 sin 3 x c4 cos 3 x and 3x 4 12 x 2 c
(d) c1 c2 x c3 sin 3 x c4 cos 3x and 5 x 4 12 x 2 c
d4y d2y
Solution (a): The homogeneous part of the given DE is 3 2 0 which is solved by putting
dx 4 dx
y e mx y me mx y m 2 e mx y m3e mx 4 mx
y m e , so e mx ( m4 3m 2 ) 0
mx
( m 4 3m 2 ) 0 , as e 0 , so m 0, 0, 3 i, 3i , thus we have two real and equal roots, i.e.
m1 m2 0 ; and two complex conjugate roots, i.e. m3 3i , m2 3i .
So complementary solution of the given DE given is as:
(0) x (0) x
yc (c1 c2 x )e e (c3 cos 3 x c4 sin 3 x) y c c1 c2 x c3 cos 3 x c4 sin 3 x .
2
Now, as the expression ‘ 108x ’ is linear 2nd degree polynomial but our differential equation is of 4th
order, so for particular solution if we must choose y p Ax 4 Bx3 Cx 2 Dx E , then
yp 4 Ax3 3 Bx 2 2Cx D y p 12 Ax 2 6 Bx 2C y
p 24 Ax 6 B y p 24 A ; on
2 2
substituting these in given DE we get, 24 A 3(12 Ax 6 Bx 2C ) 108 x , and on comparing the
2
coefficient of x , x and constant terms, we get 36 A 108 A 3 , 18 B 0 B 0 ,
24 A 6C 0 C 12 . Thus particular solution is y p 3 x 4 12 x 2 Dx E ; so the general
4 2
solution is y yc y p c1 c2 x c3 cos 3 x c4 sin 3x 3 x 12 x Dx E
y c1 (c2 D ) x c3 cos 3 x c4 sin 3 x 3 x 4 12 x 2 E , which is given by option (a).
Case 2: A function in the assumed particular solution is also a solution of the associated
homogeneous differential equation: Suppose again that g ( x ) consists of m terms of the kind
given in Table 1 and suppose further that the usual assumption for a particular solution is
y p y p1 y p2 y pn , where y p1 , y p2 , , y pn are the trial particular solution forms
corresponding to these terms. Now, if any y pi contains terms that duplicate terms in yc , then
that y pi must be multiplied by x n , where n is the smallest positive integer that eliminates that
duplication.
Linear 1st Order differential equation: In section 3.1.1, we saw that the general solution of a
linear first – order differential equation a1 ( x ) y a0 ( x ) y g ( x ) can be found by first rewriting it in
standard form
( dy dx ) P ( x ) y f ( x ) (3.53)
and assuming that P ( x ) and f ( x ) are continuous on an common interval I . Using the integrating
factor method, the general solution of Eq. 3.53 on the interval I was found to be
y c1e e P ( x ) dx f ( x)}dx . The foregoing solution has the same form given by
P ( x ) dx P ( x ) dx
{e
y yc y p . In this case yc c1e
P ( x ) dx
is a solution of the associated homogeneous equation
( dy dx ) P ( x ) y 0 (3.54)
y1 e
P ( x ) dx
is a solution of Eq. 3.54 and because the equation is linear, c1 y1 ( x ) is its general
solution. Variation of parameters consists consist of finding a particular solution of Eq. 3.53 of the
form y p u1 ( x ) y1 ( x ) . In other words, we have replaced the parameter c1 by the function u1 .
Substituting y p u1 y1 into Eq. 3.53 and using the Product Rule gives ( d dx)(u1 y1 ) P ( x )u1 y1 f ( x )
u1 ( dy1 dx ) y1 ( du1 dx ) P ( x )u1 y1 f ( x ) u1 ( dy1 dx ) P( x ) y1 y1 ( du1 dx ) f ( x)
0 because y1 is a solution of Eq. 3.57
y p y1u1 e P ( x ) dx f ( x )}dx
P ( x ) dx
{e
Linear 2nd Order Differential equation: Next we consider the case of linear second order
equation with constant coefficients
a2 y a1 y a0 y g ( x ) (3.55)
Again, we put Eq. 3.55 into the standard form by dividing by the leading coefficient a2 as,
y py qy f ( x) (3.56)
In Eq. 3.56, p and q are constant and the functions f ( x ) is continuous on some interval I . As we
have already seen earlier, there is no difficulty in obtaining the complementary solution
yc c1 y1 ( x ) c2 y2 ( x ) , the general solution of the associated homogeneous equation of Eq. 3.56,
when the coefficients are constants. Analogous to the preceding discussion, we now ask: Can the
parameters c1 and c2 can be replaced with functions and or “variable parameters,” so that
y u1 ( x ) y1 ( x ) u2 ( x) y2 ( x ) (3.57)
is a particular solution of Eq. 3.56? To answer this question we substitute Eq. 3.57 into Eq. 3.56.
Using the Product Rule to differentiate y p twice, we get y p u1 y1 y1u1 u 2 y2 y2u 2 and
y p u1 y1 y1u1 y1u1 u1 y1 u 2 y2 y2 u 2 y2 u 2 u 2 y2 ; now after substituting these derivatives and
Eq. 3.57 into Eq. 3.56 and solving, we get,
y p py p qy p f ( x ) (d dx )( y1u1 y2 u2 ) p ( y1u1 y2 u 2 ) ( y1u1 y2 u 2 ) f ( x ) (3.58)
Because we seek to determine two unknown functions u1 and u2 , reason dictates that we need two
equations. We can obtain these equations by making the further assumption that the functions u1 and
u2 satisfy y1u1 y2u2 0 y1u1 y2 u2 f ( x ) (from Eq. 3.58). By Cramer’s Rule, the solution of
the system y1u1 y2u2 0 and y1u1 y2 u2 f ( x ) can be expressed in terms of determinants:
Higher Order Equations: The method that we have just examined for non – homogeneous
second-order differential equations can be generalized to linear nth -order equations that have been put
into the standard form
y ( n ) pn 1 y ( n 1) p1 y p0 y f ( x ) (3.61)
Second Order Cauchy – Euler equation: The non – homogeneous second order Cauchy – Euler
equation has the form
ax 2 (d 2 y dx 2 ) bx (dy dx ) cy g ( x ) (3.67)
where, a( 0), b, c are constant and x 0 . The homogeneous second order Cauchy – Euler equations
associated to Eq. 3.67 is given as
ax 2 (d 2 y dx 2 ) bx ( dy dx ) cy 0 (3.68)
When we substitute y x m , the second order equation becomes, am( m 1) x m bmx m cx m 0
am ( m 1) bm c x m 0 . Thus y x m is a solution of the differential equation whenever m is
the solution of the auxiliary equation
am( m 1) bm c 0 or am 2 (b a ) m c 0 (3.69)
The discriminant of Eq. 3.69 is D (b a) 2 4ac ; and let m1 and m2 are their roots. There are three
different cases to be considered, depending on whether the roots of this quadratic equation are real
and distinct, real and equal, or complex.
Case 1: Distinct real roots: If D 0 then let m1 and m2 denotes the real roots of Eq. 3.69 such
m1 m2
that m1 m2 . Then y1 x and y2 x form a fundamental set of solutions. Hence the general
m1 m2
solution is y c1 x c2 x
Case 2: Repeated real roots: If D 0 then m1 m2 (b a ) 2a , then we obtain only one
m
solution as, y x 1 . We first write the Cauchy – Euler equation in the standard form as,
( d 2 y dx 2 ) (b ax )( dy dx ) {c (ax 2 )} y 0 and make the identifications as P ( x) b ax and
b ax dx b a ln x . Thus,
m 2m m 2 m m m m
y2 x {e ( b a )ln x x }dx x x ( b a ) x
1 1 1
dx x x ( b a ) x ( b a ) a dx x (1 x) dx x
1 1 1 1
ln x
m m
Hence the general solution for this case is y c1 x 1 c2 x 1 ln x
Case 3: Conjugate Complex roots: If the roots of Eq. 3.69 are the conjugate pair
m1 , m2 i , where and ( 0) are real, then a solution is y c1 x i c2 x i . But we
wish to write the solution in terms of real functions only. We note the identity,
x i (eln x )i e i ln x cos( ln x ) i sin( ln x ) x i cos( ln x) i sin( ln x)
x i x i 2 cos( ln x ) and x i x i 2i sin( ln x ) . From the fact that
i i
y c1 x c2 x is a solution for any values of the constants, we see, in turn, for c1 c2 1
and c1 1, c2 1 y1 x ( xi x i ) 2 x cos( ln x ) and
y2 x ( x i x i ) 2 x sin( ln x ) are the solutions.
Since W { x cos( ln x), x sin( ln x)} x 2 1 0 , 0 on the interval (0, ) , we conclude
that y1 x cos( ln x ) and y2 x sin( ln x ) constitute a fundamental set of real solutions of
the differential equation. Hence the general solution is given as,
y x c1 cos( ln x) c2 sin( ln x) .
For the particular solution, y p , we will use variation of parameter methods (discussed in Section
3.25).
Example 3.91 [ME-1998 (5 marks)]: The radial displacement in a rotating disc is governed by the
differential equation: ( d 2 u dx 2 ) (1 x )( du dx ) (u x 2 ) 8 x , where u is the displacement and x is
the radius. If u 0 at x 0 , and u 2 at x 1 . Calculate the displacement at x 1 2 .
Solution: The given differential equation can be written as: x 2 ( d 2 u dx 2 ) x ( du dx) u 8 x 3 …(i)
2 2 2
whose associated homogeneous part is x ( d u dx ) x (du dx ) u 0 …(ii), which is in the form of
Example 3.92 [ME-1998 (2 marks)]: The general solution of the differential equation
x 2 ( d 2 y dx 2 ) x( dy dx) y 0 is (where, A and B are constants)
(a) Ax Bx 2 (b) Ax B log x (c) Ax Bx 2 log x (d) Ax Bx log x
Solution (d): As the given equation is in the form of homogeneous Cauchy – Euler equation, so
putting y x m y mx ( m 1) y m(m 1) x ( m 2) in the given equation we get,
x 2 m(m 1) x ( m 2) xmx ( m1) x m 0 x m ( m 2 2m 1) 0 as m
x 0 so ( m 2 2m 1) 0
(m 1) 2 0 m 1 . As we have reapeated real roots so the genereal solution of the given
homogeneous Euler-Cauchy equation is y c1 x c2 x ln x .
Example 3.94 [XE-2013 (2 mark)]: The general solution of the differential equation
x 3 ( d 3 y dx 3 ) x 2 ( d 2 y dx 2 ) x ( dy dx ) y 0 , x 0 is (where, k1 ( 3 2) x , k2 ( 3 2) log e x ).
x x 2 1 2
(a) c1e e c2 cos(k1 ) c3 sin(k1 ) (b) c1 x x c2 cos(k2 ) c3 sin(k2 )
x x 2 12
(c) c1e e c2 cos(k1 ) c3 sin(k1 ) (d) c1 x x c2 cos(k2 ) c3 sin(k2 )
Solution (d): As the given differential equation is of the form of Cauchy-Euler equation so putting
y x m y mx ( m1) y m( m 1) x ( m 2) y m( m 1)(m 2) x ( m 3) in the given equation we
m 2
get x m(m 1)(m 2) m(m 1) m 1 0 (m 1)(m m 1) 0 as x m 0 . So the roots of
auxiliary equation are: m 1, (1 3i ) 2 . So the general solution of the given differential equation is
y c1 x x1 2 [c2 cos{( 3 2) x} c3 sin{( 3 2) x}]
Example 3.95 [XE-2014 (1 mark)]: Which of the following is a solution of the differential equation
x 2 y xy y 4 sin(ln x ) , x 0 ?
(a) y 2 x sin(ln x) (b) y 2 x sin(ln x) (c) y 2 ln x cos(ln x) (d) y 2 ln x cos(ln x)
Solution: The associated homogeneous part of the given differential equation is x 2 y xy y 0
which is in the form of Euler-Cauchy equation so putting
m ( m 1) ( m 2)
y x y mx y m(m 1) x in the given equation we get
x m m(m 1) m 1 0 (m2 1) 0 as x m 0 . So the roots of auxiliary equation are: m 0 i .
0
So the complementary solution yc x c1 cos(ln x) c2 sin(ln x ) c1 cos(ln x ) c2 sin(ln x) . As we
have duplication of y c with the function g ( x ) 4 sin(ln x) so the choice of our particular solution is
y p A ln x cos(ln x) B ln x sin(ln x ) . Thus options (a) and (b) are not correct; now we have to
choose between options (c) and (d). By checking these two option we found option (c) is correct as:
y p 2 ln x cos(ln x) yp (2 x) cos(ln x ) ln x sin(ln x )
yp 4 sin(ln x) 2 ln x cos(ln x ) 2 cos(ln x ) 2 ln x sin(ln x) x 2 . Now putting these into the
given equation we get LHS 4 sin(ln x ) 2 ln x cos(ln x ) 2 cos(ln x ) 2 ln x sin(ln x )
2 cos(ln x ) 2 ln x sin(ln x ) 2 ln x cos(ln x ) 4 sin(ln x ) RHS .
2 d2y dy
Example 3.97 [PE-2016 (2 marks)]: For the differential equation x 2
2x 2 y 0 , the
dx dx
general solution is
(a) y c1 x c2 e x (b) y c1 sin x c2 cos x (c) y c1e x c2 e x (d) y c1 x 2 c2 x
Solution (d): The given differential equation is of 2nd order homogeneous Cauchy – Euler equation,
whose solution is of the form y x m y mx ( m 1) y m(m 1) x ( m 2) ; so substituting these in
m
the given DE, we get x m {m( m 1) 2m 2} 0 {m( m 1) 2m 2} 0 , as x 0 .
Thus {m( m 1) 2m 2} 0 m 2 3m 2 0 m 1, 2 , thus we have real and distinct roots; so
the general solution of given DE is y c1 x 2 c2 x1 .
[Similar question was also asked in EE-2014 (1 mark)]
Reduction to Constant coefficients: The similarities between the forms of solutions of Cauchy-
Euler equations and solutions of linear equations with constant coefficients are not just a coincidence.
For example, when the roots of the auxiliary equations for ay by cy 0 and
m1x m2 x
ax 2 y bxy cy 0 are distinct and real, the respective general solutions are y c1e c2 e and
m1 m2
y c1 x c2 x , x 0. In view of the identity e ln x x , x 0 , the second solution can be
m ln x m ln x mt m t
expressed in the same form as the first solution y c1e 1 c2 e 2 c1e 1 c2 e 2 , where t ln x .
This last result illustrates the fact that any Cauchy-Euler equation can always be rewritten as a linear
differential equation with constant coefficients by means of the substitution x et . The idea is to
solve the new differential equation in terms of the variable t , using the methods of the previous
sections, and, once the general solution is obtained, re-substitute t ln x .
Solution for < : In the preceding discussion we have solved Cauchy-Euler equations for x 0 .
One way of solving a Cauchy-Euler equation for x 0 is to change the independent variable by
means of the substitution t x (which implies t 0 ) and using the Chain Rule:
dy dy dt dy d2y d dy dt d 2 y
and .
dx dt dx dt dx 2 dt dt dx dt 2
Example 3.99 [TF-2007 (2 marks)]: Using exp( x 2 ) as an integrating factor, the solution of the first
order differential equation y 2 xy 1 in terms of the error function [erf ( x )] and a constant of
integration c , is given by
(a) y c ( 2)erf ( x ) exp( x 2 ) (b) y ( 2)erf ( x ) c exp( x 2 )
(c) y ( 2)erf ( x ) c exp( x 2 ) (d) y c ( 2)erf ( x) exp( x 2 )
2
Solution (c): Multiplying both sides of the given differential equation with e x , we get
2 2 2 2 2 2 2
( dy dx )e x 2 xye x e x ( d dx )( ye x ) e x d ( ye x ) e x dx ; integrating both sides we
x2 2 2 2 2 2
get d ( ye ) e x dx ye x e x dx c y { e x dx c}e x . Now error function is
x 2 x2
defined as erf ( x ) (2 ) e t dt . Thus e dx ( 2)erf ( x ) . So
0
2
x
y ( 2)erf ( x ) c e .
Statement for Linked Answer Question 3.100 and 3.101: Let y ( x ) n 0 an x n be a solution of
the differential equation ( d 2 y dx 2 ) xy 0
Example 3.101 [XE-2007 (2 marks)]: The solution of the differential equation given above
satisfying y (0) 1 and y (0) 0 is
(a) y ( x ) 1 1 (2 3) x 2 1 (2 3 5 6) x 4 1 (2 3 5 6 8 9) x 6
(b) y ( x ) 1 1 (2 3) x 2 1 (2 3 5 6) x 4 1 (2 3 5 6 8 9) x 6
(c) y ( x ) 1 1 (2 3) x 3 1 (2 3 5 6) x 6 1 (2 3 5 6 8 9) x 9
(d) y ( x ) 1 1 (2 3) x 3 1 (2 3 5 6) x 6 1 (2 3 5 6 8 9) x 9
Solution (d): From the previous problem we have
n n
y ( x ) n 0 an x a0 a1 x an x y (0) a0 1
y ( x ) n 1 nan x n 1 a1 2a2 x nan x n 1 y(0) a1 0
As from the recursion relation we have a2 0 and an 2 an 1 (n 1)(n 2)
a3 1 (2 3) a0 1 (2 3) , a4 1 (3 4) a1 0 , a5 1 (4 5) a2 0 ,
a6 1 (5 6) a3 1 (2 3 5 6) , a7 1 (6 7) a4 0 , a8 1 (7 8) a5 0 ,
a9 1 (8 9) a6 1 (2 3 5 6 8 9) , and so on. Thus substituting the obtained result in (i) we get
y ( x ) 1 1 (2 3) x 3 1 (2 3 5 6) x 6 1 (2 3 5 6 8 9) x 9 .
Exercise: 3.2
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. The solution of the IVP y 9 y 0 , y (0) 2 , y (0) 1 is
(a) y ( x ) (7 6)e 3 x (5 6)e 3 x (b) y ( x ) (7 6)e3 x (5 6)e 3 x
(c) y ( x ) (6 7)e3 x (6 5)e 3 x (d) y ( x ) (6 7)e 3 x (6 5)e 3 x
2. One of the solutions of y 9 y 0 is
(a) e x (b) e 2 x (c) e3x (d) e 4 x
3. The solution of the IVP y 11 y 24 y 0 , y (0) 0 , y (0) 7 is
(a) y ( x ) (7 5)(e 8 x e 3 x ) (b) y ( x ) (5 7)(e 8 x e 3 x )
(c) y ( x ) (5 7)(e 8 x e 3 x ) (d) y ( x ) (7 5)(e 8 x e 3 x )
4. The solution of the IVP 4 y 5 y 0 , y ( 2) 0 , y (2) 7 is
(a) y ( x ) (28 5){1 e (5 x 4) (5 2) } (b) y ( x ) (28 5){1 e (5 x 4) (5 2) }
(c) y ( x ) (28 5){1 e (5 x 4) (5 2) } (d) y ( x ) (28 5){1 e (5 x 4) (5 2) }
5. The solution of the IVP y 4 y 9 y 0 , y (0) 0 , y (0) 8 is
t 2t
(a) y ( x ) (8 5)e sin( 5t ) (b) y ( x ) (8 5)e sin( 5t )
(c) y ( x ) (8 5)e 2 t sin( 5t ) (d) y ( x ) (8 5)et sin( 5t )
6. The solution of IVP y 8 y 17 y 0 , y (0) 4 , y (0) 1 is
(a) y ( x ) 4e 4 x cos x 15e 4 x sin x (b) y ( x ) 4e 4 x cos x 15e 4 x sin x
(c) y ( x ) 4e 4t cos x 15e 4 x sin x (d) y ( x ) 4e 4t cos x 15e 4 x sin x
7. The solution of IVP y 16 y 0 , y( 2) 10 , y( 2) 3 is
(a) y( x) 10 cos(4 x) (3 4) sin(4 x) (b) y( x) 10 cos(4 x) (3 4) sin(4 x)
(c) y( x) 10 cos(4 x) (3 4) sin(4 x) (d) y( x) 10 cos(4 x) (3 4)sin(4 x)
8. The solution of IVP y 4 y 4 y 0 , y (0) 12 , y (0) 3 is
(a) y ( x ) 12e 2 x 27(1 x)e 2 x (b) y ( x ) 12e 2 x 27 x 2 e 2 x
(c) y ( x ) 12e 2 x 27 xe 2 x (d) y ( x ) 12e 2 x 27 x3 e 2 x
9. The solution of IVP y 14 y 49 y 0 , y ( 4) 1 , y (4) 5 is
(a) y ( x ) 9e 7( x 4) 2 x 2 e 7( x 4) (b) y ( x ) 9e 7( x 4) 2 xe 7( x 4)
(c) y ( x ) 9e 7( x 4) 2 xe7( x 4) (d) y ( x ) 9e 7( x 4) 2 xe 7( x 4)
10. The general solution of 2 x 2 y xy 3 y 0 , given that y1 ( x ) 1 x is a solution, where c1 is a
constant, is
(a) y ( x ) c1 (1 x ) c2 x 3 2 (b) y ( x ) c1 (1 x ) c2 x1 2
(c) y ( x ) c1 (1 x) c2 x 5 2 (d) y ( x ) c1 (1 x ) c2 x
11. The general solution of x 2 y 2 xy 2 y 0 , given that y1 ( x ) x is a solution, where c1 is a
constant, is
(a) y ( x ) c1 x c2 (1 x ) (b) y ( x ) c1 x c2 (1 x 2 )
(c) y ( x ) c1 x c2 ( x 2 ) (d) y ( x ) (c1 c2 x ) x1 2
12. The Wronskian of two solution to the following DE t 4 y 2t 3 y t 8 y 0 , where c is a constant,
is
(a) c (b) ct (c) ct 2 (d) c(1 t )
5t
13. The particular solution of the DE y 4 y 12 y 3e is
5t 5t 5t 5t
(a) y p (3 7)e (b) y p (3 7)e (c) y p (7 3)e (d) y p (7 3)e
14. The particular solution of the DE y 4 y 12 y sin(2t ) is
(a) y p (1 40) cos(2t ) (b) y p (1 20) sin(2t )
(c) y p (1 40) cos(2t ) (1 20) sin(2t ) (d) None of these
15. The particular solution of the DE y 4 y 12 y 2t 3 t 3 is
3 2 3 2
(a) y p (1 6)t (1 6)t (1 9)t (5 27) (b) y p (1 6)t (1 6)t (1 9)t (5 27)
3 2 3 2
(c) y p (1 6)t (1 6)t (1 9)t (5 27) (d) y p (1 6)t (1 6)t (1 9)t (5 27)
16. The particular solution of the DE y 4 y 12 y te 4 t is
4t 4t
(a) y p (1 36)(3t 1) e (b) y p (1 36)(3t 1)e
4t 4t
(c) y p (1 36)e (d) y p (1 36)e
17. The particular solution of the DE y p ( x ) y q ( x) y 16e 7 x sin(10 x ) , is
7x 7x
(a) y p Ae cos(10 x ) (b) y p Ae sin(10 x) where A, B are
7x constants
(c) y p A cos(10 x) B sin(10 x ) (d) y p e { A cos(10 x ) B sin(10 x )}
multiplying all results, we have I n {n(n 1)(n 2) 2 1} s n I 0 ( n ! s n ) I 0 .
b b
Since I 0 lim e st dt lim ( e st s ) lim{( e sb 1) s} 1 s
b 0 b 0 b
I n n ! s n 1 {t n } n ! s n 1 , n 1, 2, , n
{e at } 1 ( s a ) ; from Eq. 3.71, we have,
b b b
{e at } lim e st e at dt lim e ( s a ) t dt lim e ( s a ) t ( s a )
b 0 b 0 b
0
Example 3.102 [ME-1999 (1 mark)]: Laplace transform of ( a bt ) 2 where, a and b are constants
is given by
(a) ( a bs ) 2 (b) 1 ( a bs) 2 (c) ( a 2 s ) (2ab s 2 ) (2b 2 s 3 ) (d) ( a 2 s ) (2ab s 2 ) (b 2 s 3 )
Solution (c): f (t ) ( a bt ) 2 a 2 2abt b 2t 2 ; so
{ f (t )} {a 2 2abt b 2t 2 } a 2 {1} 2ab{t } b 2 {t 2 } ( a 2 s) (2ab s 2 ) (2b 2 s 3 )
Example 3.103 [ME-2000 (2 marks)]: The Laplace transform of the function sin 2 2t is
(a) {1 2s} [ s {2( s 2 16)}] (b) s ( s 2 16) (c) (1 s ) {s ( s 2 4)} (d) s ( s 2 4)
Solution (a): f (t ) sin 2 2t (1 cos 4t ) 2 ; so
{ f (t )} (1 2) 1 cos 4t (1 2)({1} {cos 4t}) { f (t )} {1 2 s} [ s {2( s 2 16)}]
st 1 2 u 2 s 2u 2 u2 1 u2 1 u2
0 e t dt e
0 u s
du e
s 0
du e
s
du
s
s
. As e du is a
Example 3.108 [EC-1996 (2 marks)]: The inverse Laplace transform of the function
( s 5) {( s 1)(s 3)} is
t t
1. So the limit cannot be determined.
Example 3.111 [EC-2003 (1 mark)]: The Laplace transform of i(t ) is given by I ( s ) 2 {s (1 s)} .
As t , the value of i (t ) tends to
(a) 0 (b) 1 (c) 2 (d)
2 1 1 1 1 1 1 1
Solution (c): As i (t ) 1 2 2 2 2 1 2e
t
s(1 s) s 1 s s 1 s
i (t ) 2(1 e t ) lim i (t ) lim 2(1 e t ) 2
t t
Statement for Linked Answer Questions 3.112 and 3.113: Let F ( s ) ( s 10) {( s 2)( s 20)}
Example 3.112 [AE-2007 (2 marks)]: The partial fraction expansion of F ( s ) is
1 1 5 2 2 20 49 59
(a) (b) (c) (d)
s 2 s 20 s 2 s 20 s 2 s 20 s 2 s 20
( s 10) A B s ( A B ) (20 A 2 B)
Solution (d): F ( s) A B 1 and
( s 2)( s 20) s 2 s 20 ( s 2)( s 20)
20 A 2 B 10 ; solving these two equations we get, A4 9 and B 5 9; thus
F ( s ) {(4 9) ( s 2)} {(5 9) ( s 20)} .
Example 3.114 [ME-2009 (1 mark), ME-2012, PI-2012 (2 marks)]: The inverse Laplace transform
of 1 ( s 2 s ) is
(a) 1 e t (b) 1 e t (c) 1 e t (d) 1 e t
1 1 1 1 1 1
Solution (c): F ( s) 2 1{F ( s )} 1 1 1 e
t
s s s ( s 1) s s 1
s s 1
Now for solving linear initial – value problems in which the differential equation has constant
coefficients; such a differential equation is simply a linear combination of terms y , y , , y ( n ) .
an ( d n y dt n ) an 1 ( d n 1 y dt n 1 ) a0 y g (t ) , y (0) y0 , y (0) y1 , , y ( n 1) (0) yn 1 , where
the an , an 1 , , a1 , a0 and y0 , y1 , , yn 1 are constants. By the linearity property the Laplace
transform of this linear combination is a linear combination of Laplace transform:
an (d n y dt n ) an1(d n 1 y dt n1 ) a0 { y} {g (t )} ; from Eq. 3.75 it becomes
an s nY ( s ) s n 1 y (0) y ( n 1) (0) an 1 s n 1Y ( s ) s n 2 y (0) y ( n 2) (0)
(3.76)
a0Y ( s) G ( s)
where, { y (t )} Y ( s ) and {g (t )} G ( s) . In other words, the Laplace transform of a linear
equation with constant coefficient becomes an algebraic equation in Y ( s) . If we solve the general
transformed Eq. 3.76 for the symbol Y ( s) , we first obtain P ( s )Y ( s ) Q( s ) G ( s ) and then write
Y ( s ) {Q( s) P ( s )} {G( s) P ( s )} (3.77)
n n 1
where P ( s ) an s an 1s a0 , Q ( s ) is a polynomial in s of degree less than or equal to
consisting of the various products of the coefficient an , an 1 , , a1 , a0 and the prescribed initial
conditions y0 , y1 , , yn 1 , and G ( s ) is the Laplace transform of g (t ) . Typically, we put the two
terms in Eq. 3.77 over the least common denominator and then decompose the expression into two or
more partial fractions. Finally, the solution y (t ) of the original initial-value problem is
y (t ) 1Y ( s ) , where the inverse transform is done term by term.
Initial Value Theorem: The initial value theorem determines the value of the time function,
f (t ) , when t 0 without finding the inverse Laplace transform. Now, as s , we have
st
e 0 thus lim { f (t )} lim ( d dt ) f (t ) lim ( d dt ) f (t ) e st dt 0 . From Eq. 3.77
s 0 s0 s0 0
we have lim { f (t )} lim sF ( s) f (0) . Thus equating the two results we have
s s
lim sF ( s) f (0) 0 lim{sF (s )} f (0) lim{sF (s )} lim f (t ) . Thus the initial value of
s s s t 0
Eq. 3.77 we have lim { f (t )} lim sF ( s ) f (0) . Thus equating the two results we have
s 0 s 0
lim sF (s ) f (0) f () f (0) lim{sF (s )} f () lim{sF (s )} lim f (t ) . Thus the final
s 0 s 0 s 0 t
value of the function f (t ) is lim f (t ) lim{sF ( s)} [This point was asked in EE-2002 (1
t s 0
mark)].
Transfer Function: A transfer function is the ratio of the output of a system to the input of a
system, in the Laplace domain considering its initial conditions and equilibrium point to be zero.
If we have an input function of X ( s ) and an output function of Y ( s) , we define the transfer
function H ( s) to be: H ( s ) Y ( s ) X ( s ) .
Example 3.116 [ME-2002 (5 marks)]: Using Laplace transform, solve ( d 2 y dt 2 ) 4 y 12t . Given
that y 0 and dy dt 9 at t 0 .
Solution: Taking Laplace transform on both sides of the given differential equation we get
{ y } 4{ y} 12{t} . Using equation 3.80 and given conditions, we have
{ y } s 2Y ( s ) sy (0) y (0) { y } s 2Y ( s ) 9 ; { y} Y ( s ) ; {t} 1 s 2 . Now substituting
12 9
these in { y } 4{ y} 12{t} , we get s 2Y ( s ) 9 4Y ( s ) 12 s 2 Y ( s ) 2 2 2 .
s ( s 4) s 4
Using partial fractions for the first expression of Y ( s) , we get
3 3 9 3 6 1 1 3 6
Y (s) 2
2
2
2
2
y (t ) {Y ( s)} 2
2
s s 4 ( s 4) s ( s 4) s ( s 4)
1 2 1 2 2
y (t ) 3 (1 s ) 3 {2 ( s 2 )} 3t 3sin 2t
[Similar question was also asked in ME-1997 (5 mark)]
Example 3.117 [EC-2006 (2 marks)]: The unit-step response of a system from rest is given by
c (t ) 1 e 2 t for t 0 . The transfer function of the system is:
(a) 1 (1 2 s ) (b) 2 (2 s ) (c) 1 (2 s ) (d) 2 s (1 2 s )
Solution (b): As the input of the system is c(0) 1 and output of the system is c (t ) 1 e 2 t . So we
have input function X ( s) {c(0)} {1} 1 s ; also the output function is given as:
Y ( s ) {c (t )} {1 e } {1} {e 2 t } (1 s ) {1 ( s 2)} 2 {s ( s 2)} . Thus transfer function
2 t
H ( s ) Y ( s) X ( s) [2 {s ( s 2)}] s 2 ( s 2) .
Example 3.118 [EC-2006 (2 marks)]: Consider the function f (t ) having Laplace transform
F ( s ) 0 ( s 2 02 ) , Re[ s ] 0 . The final value of f (t ) would be:
(a) 0 (b) 1 (c) 1 f () 1 (d)
1 1 2 2
Solution (c): f (t ) {F (s )} {0 (s 0 )} sin(0t ) . Applying final value theorem, the
final value of f (t ) lim sin(0t ) ; as t , sin(0 t ) lies between –1 to 1 . So 1 f () 1 .
t
Example 3.121 [ME-2013, PI-2013 (2 marks)]: The function f (t ) satisfies the differential equation
( d 2 f dt 2 ) f 0 and the auxiliary conditions, f (0) 0 , ( df dt )t 0 4 . The Laplace transform of
f (t ) is given by
(a) 2 ( s 1) (b) 4 ( s 1) (c) 4 ( s 2 1) (d) 2 ( s 4 1)
Solution (c): Taking Laplace transform on both sides of the given differential equation we get
{ f } { f } {0} . Using equation 3.80 and given conditions, we have
{ f } s 2 F ( s ) sf (0) f (0) { f } s 2 F ( s) 4 ; { f } F ( s ) ; {0} 0 . Now substituting
these in { f } { f } {0} , we get s 2 F ( s) 4 F ( s ) 0 F ( s ) 4 ( s 2 1)
[Similar question was also asked in TF-2008 (2 marks)]
Example 3.122 [AE-2014 (2 marks)]: The Laplace transform L{u (t )} U ( s) , for the solution u (t )
of the problem ( d 2 u dt 2 ) 2( du dt ) u 1 , t 0 with initial conditions u (0) 0 , ( du dt )t 0 5 is
given by:
(a) 6 ( s 1) 2 (b) (5s 1) {s ( s 1) 2 } (c) (1 5s ) {s ( s 1) 2 } (d) (5s 2 1) {s ( s 1) 2 }
Solution (c): Taking Laplace transform on both sides of the given differential equation we get
{u} 2{u} {u} {1} . Using equation 3.80 and given conditions, we have
{u} s 2U ( s ) su (0) u(0) {u } s 2U ( s) 5 ;
{u} sU ( s ) u (0) {u } sU ( s ) ;
{u} U ( s) ; {1} 1 s . Now substituting these in {u} 2{u} {u} {1} , we get
s 2U ( s ) 5 2 sU ( s) U ( s ) 1 s U ( s ) (1 5s ) {s ( s 2 2s 1)} (1 5s) {s ( s 1) 2 } .
Example 3.123 [EE-2014 (1 mark)]: Let X ( s) (3s 5) {s 2 10 s 21} be the Laplace Transform
of a signal x (t ) . Then x(0 ) is
(a) 0 (b) 3 (c) 5 (d) 21
Solution (b): We have to find the initial value of f (t ) whose Laplace transform is X ( s ) , thus
lim f (t ) lim{sX ( s )} lim[ s (3s 5) {s 2 10 s 21}] lim[(3 5 s ) (1 10 s 21 s 2 )] 3 1 3
t 0 s s s
Example 3.124 [IN-2016 (2 marks)]: The relationship between the force f (t ) and the displacement
x (t ) of a spring-mass system (with a mass M , viscous damping D and spring constant K ) is
d 2 x (t ) dx(t )
M 2
D Kx (t ) f (t ) . X ( s ) and F ( s) are the Laplace transform of x (t ) and f (t )
dt dt
respectively. With M 0.1 , D 2 , K 10 in appropriate units, the transfer function
G ( s ) X ( s ) F ( s ) is
(a) 10 ( s 2 20s 100) (b) ( s 2 20 s 100) (c) 10s 2 ( s 2 20 s 100) (d) s ( s 2 20 s 100)
d 2 x (t ) dx(t )
Solution (a): Our DE is 0.1 2
2 10 x (t ) f (t ) …(i); taking Laplace transform on both
dt dt
sides of (i), we get 0.1{s 2 X ( s ) sx (0) x(0)} 2{sX ( s ) x (0)} 10 X ( s ) F ( s) …(ii). As
x(0) x(0) 0 , so from (ii), we have 0.1s 2 X ( s ) 2 sX ( s) 10 X ( s) F ( s )
G ( s) X ( s ) F ( s ) 1 (0.1s 2 2 s 10) 10 ( s 2 20 s 100) .
Example 3.126 [EE-2016 (1 mark)]: The solution of the differential equation, for t 0 ,
y (t ) 2 y (t ) y (t ) 0 with initial conditions y (0) 0 and y (0) 1 , is { u (t ) denotes the unit step
function}
(a) te t u (t ) (b) (e t te t )u (t ) (c) ( e t te t )u (t ) (d) e t u (t )
Solution (a): We have homogeneous linear equation with constant coefficients, so its solution will be
of the form y e mt y me mt y m 2 e mt , putting all these in given DE, we get
mt
e mt ( m 2 2m 1) 0 (m 2 2m 1) 0 , as e 0 ; so ( m 2 2m 1) 0 ( m 1) 2 0 m 1
Thus we have real and equal roots and so the general solution of given DE is
y (c1 c2 t )e mt (c1 c2 t )e t y c2 e t (c1 c2 t )e t
Now y (0) 0 0 c1 ; and y (0) 1 1 c2 e 0 (0 c2 (0))e 0 c2 1
0, t 0
Hence y te t ; as u (t ) , but given DE is valid for t 0 , so required solution will be
1, t 0
y te t u (t ) . Thus option (a) is correct.
Unit Impulse function: Mechanical system are often acted on by an external force or electromotive
force of large magnitude that acts only for a very short period of time. For e.g., a vibrating airplane
wing could be struck by lightning; the graph of the piecewise-defined function
1 k , a t a k
a (t a ) (3.78)
0, otherwise
where, a 0 and k is a small positive number, for the figure shown in
Fig. 3.3, could serve as a model for such a force. For a small value of k ,
a (t a ) is essentially a constant function of large magnitude that is ‘on’ Figure 3.3: Graph of ( − )
for just a very short period of time, around a . The function a (t a ) is called a unit impulse,
because it possesses the integration property 0 a (t a)dt 1 .
a ak
Proof: 0 a (t t0 )dt 0 a (t a )dt a a (t a)dt a k a (t a)dt , the first and the third
integral becomes zero as a (t a) 0 for all t a and t a k ; and a (t a ) 1 k for all
a k
a t a k . So 0 a (t a)dt (1 k )a dt ( k k ) 1 .
It is to be noted that the derivative of the unit step function u (t a) is called the unit impulse
function a (t a ) , i.e., a (t a) ( d dt )u (t a ) .
The unit impulse function can be written in terms of unit step function as:
a (t a ) {1 (2 a )}u (t a ) u t (a k ) .
2
Example 3.128 [IN-2011 (1 mark)]: The integral (1 2 ) t 2 e t 2 (1 2t ) dt is equal to
2 s
Example 3.136 [TF-2014 (1 mark)]: The inverse Laplace transform of 2
2
is
s 2 s s 4s 4
(a) 1 2te 2t (b) 1 e 2 t 2te 2 t (c) 1 2e 2 t 2te 2t (d) e 2t 2te 2t
2 s 2 s22 1 1 s2 2
Solution: F ( s ) 2 2 2
2
2
s 2s s 4 s 4 s ( s 2) ( s 2) s s 2 ( s 2) ( s 2)
1 1 1 2 1 2 1 1 1 2
F (s) 2
{F ( s )}
2 2
s s2 ( s 2) ( s 2) s ( s 2) s ( s 2)
1 1 1
{F ( s )}
2
2
s ( s 2)
1 1
s
1
2
1
2
( s 2)
1 2te
2 t
Inverse form of Eq. 3.80: To compute the inverse of F ( s a) , we must recognize F ( s ) , find
f (t ) by taking the inverse Laplace transform of F ( s ) , and then multiply f (t ) by the exponential
function e at . This procedure can be summarized symbolically in the following manner:
1 F ( s a ) e at f (t ) , where, f (t ) 1{F ( s )} (3.81)
s 2 s 5 ( s 1) 2 s 2s 5
F ( s 1) ( s 1) {( s 1) 2 22 } F ( s) s ( s 2 22 ) . Thus using Eq. 3.84, we have
f (t ) 2 1
F (s 1) 2e t
g (t ) , where 1
g (t ) {F ( s )} 1
s ( s 2 ) cos(2t ) .
2 2
So
f (t ) 2e t cos 2t . Now lim f (t ) lim f (0 h) 2e h cos 2h 2 and
t 0 h 0
Inverse form of Eq. 3.90: If f (t ) 1{F ( s )} , the inverse form of Eq. 3.83, a 0 is given as,
1{e a s F ( s )} f (t a )u (t a) (3.84)
Laplace Transform of Unit Step Function
a
{u (t a )} e st u (t a )dt e st u (t a ) dt e st u (t a ) dt 0 e st dt e st ( s )
0 0 a a
a
as as as
{u (t a)} 0 (e s) e s {u (t a )} e s [This point was asked in CE-1998,
IN-2010 (1 mark), ME-2004 (2 marks)]
Laplace Transform of Unit Impulse Function
{ a (t a )} (1 k ) u (t a ) u t ( a k ) {1 ( ks )}{e as e ( a k ) s } {e as ( ks )}(1 e ks )
Laplace Transform of Dirac Delta Function
{ (t a)} {lim a (t a)} lim { a (t a)} lim e as (1 e ks ) ( ks ) ;
k 0 k 0 k 0
now by using
as
L’Hospital rule, we have { (t a)} e lim{se ks s} e as 1 e as .
k 0
When a 0 { (t )} 1 .
Laplace Transform of Unit Ramp Function
{r (t a )} (t a)u (t a ) f (t a ) u (t a ) , where f (t ) t f (t a ) (t a ) ; now
using the second shifting theorem, we have {r (t a )} e as { f (t )} e as {t} e as s 2 , as
{t} 1! s11 1 s 2 . Thus {r (t a )} e as s 2 [This point was asked in EC-1994 (1 mark)].
2, 0 t 1
Example 3.138: Write the following function 2
f (t ) t 2 , 1 t 2
using unit step functions and find its transform.
cos t , t 2
Solution: In terms of unit step function the given function can be written as,
f (t ) 2 1 (t 1) (t 2 2) u (t 1) u (t 2) (cos t )u (t 2)
To apply Eq. 3.91, we must write each term in f (t ) in the form f (t a )u (t a) . Thus,
2 1 (t 1) 2(1 e s ) s
1
2 1 1 1 1 1
t 2 (t 1) (t 1) 2 (t 1) (t 1) 3 2 e s
2 2 s s 2s
2 2 s 2
1
2
1
2
2
t 2 (t 2) (t 2) 2 (t 2)
8
(t 2)
1
3
s 2s
2
8s
e
2(1 e s ) 1 1 1 s 1 2 s 2
Hence, f (t ) 3 2 e 3 2 e
s s s 2s s 2s 8s
Alternative form of Eq. 3.83: We are frequently confronted with the problem of finding the
Laplace transform of a product of a function g and a unit step function (t a ) where the
function g lacks the precise shifted form in Eq. 3.83. To find the Laplace transform of
g (t ) (t a) , it is possible to find up g (t ) into the required form by algebraic manipulations, just
as we did it in previous example for 2nd and 3rd term. But since these manipulations are time
consuming and often not obvious, it is simpler to devise an alternative version of Eq. 3.83. Using
the definition of (t a ) and the substitution u t a , we obtain
g (t ) (t a ) e st g (t ) dt e s ( u a ) g (u a ) du e a s e s u g (u a )du
a a a
as s t a s
g (t ) (t a ) e a e g (t a ) dt e g (t a ) (3.85)
Example 3.144 [EC-2004 (2 marks)]: A system is described by the following differential equation
( d 2 y dt 2 ) 3( dy dt ) 2 y x (t ) is initially at rest. For input x (t ) 2u (t ) , the output y (t ) is
(a) (1 2e t e 2 t )u (t ) (b) (1 2e t e 2 t )u (t )
(c) (0.5 e t 1.5e 2 t )u (t ) (d) (0.5 2e t 2e 2t )u (t )
Solution (a): Taking Laplace transform on both sides of the given differential equation we get
{ y } 3{ y } 2{ y} 2{u (t )} ; using Eq. 3.75 and the condition: system is initially at rest, i.e.,
y (0) 0 and y (0) 0 , we have { y } s 2Y ( s ) sy (0) y (0) s 2Y ( s) ;
{ y } sY ( s ) y (0) sY ( s) ; { y} Y ( s ) ; since {u (t )} 1 s . Thus substituting all these in
{ y } 3{ y } 2{ y} 2{u (t )} we get s 2Y ( s) 3sY ( s) 2Y ( s ) 2 s
Y ( s ) 2 {s ( s 2 3s 2)} 2 {s( s 2)( s 1)} ; from partial fractions we have
Y ( s ) {1 s} {2 ( s 1)} {1 ( s 2)} Now taking inverse Laplace transform on both sides we get
y (t ) 1{Y ( s )} 1{1 s} 21{1 ( s 1)} 1{1 ( s 2)} 1 2e t e 2 t . Now as
0, t 0
u (t ) so y (t ) can be represented as y (t ) (1 2e t e 2 t )u (t ) .
1, t 0
Example 3.145 [EE-2008 (1 mark)]: A function y (t ) satisfies the following differential equation
( d dt ) y (t ) y (t ) (t ) , where (t ) is the delta function. Assuming zero initial condition, and
denoting the unit step function by u (t ) , y (t ) can be of the form
(a) e t (b) e t (c) et u (t ) (d) e t u (t )
Solution (d): We have given y (0) 0 . Now taking Laplace transform on both sides of the given
differential equation we get { y } { y} { (t )} ; using Eq. 3.75 and given conditions, we have
{ y } sY ( s ) y (0) sY ( s) ; { y} Y ( s ) ; since { (t )} e 0 s 1 . Thus substituting all these in
{ y } { y} { (t )} we get sY ( s ) Y ( s ) 1 Y ( s ) 1 ( s 1) . Now taking inverse Laplace
transform on both sides we get y (t ) 1{Y ( s )} 1{1 ( s 1)} e t . As for t 0 , y 0 and for
t 0 , y (t ) e t ; so y (t ) e t u (t ) .
Solution (b): Taking Laplace transform on both sides of the given differential equation we get
{ y } 5{ y} 6{ y} {x (t )} ; using Eq. 3.75 and given conditions, we have
{ y } s 2Y ( s ) sy (0) y (0) s 2Y ( s) ; { y } sY ( s ) y (0) sY ( s) ; { y} Y ( s ) ; since
2 s 2 s
x (t ) {u (t ) u (t 2)} { x(t )} {u (t ) u (t 2)} (1 s ) (e s) (1 e ) s . Thus
substituting all these in the given differential equation after taking Laplace, we get
s 2Y ( s) 5sY ( s ) 6Y ( s ) (1 e 2 s ) s
Y ( s ) (1 e 2 s ) {s( s 2 5s 6)} (1 e 2 s ) {s( s 2)( s 3)} .
[Similar question was also asked in AE-2012 (1 mark)]
Derivatives of a Transform: The Laplace transform of the product of a function f (t ) with t can
be found by differentiating the Laplace transform of f (t ) . Let us assume F ( s) { f (t )} exist and
that it is possible to interchange the order of differentiation and integration. Then
( d ds ) F ( s ) ( d ds ) e st f (t ) dt ( s ){e st f (t )}dt e st t f (t ) dt {t f (t )}
0 0 0
Example 3.150 [AE-2008 (2 marks)]: Let Y ( s) denote the Laplace transform L y (t ) of the
function y (t ) cosh(at ) sin( at ) . Then
(a) L (dy dt ) ( dY ds ), L t y (t ) sY ( s ) (b) L ( dy dt ) sY ( s ), L t y (t ) ( dY ds )
(c) L (dy dt ) (dY ds ), L t y (t ) Y ( s 1)
(d) L ( dy dt ) sY ( s ), L t y (t ) e asY ( s )
Solution (b): Using Eq. 3.75, we have { y (t )} sY ( s) f (0) , as y (0) 0 so { y (t )} sY ( s ) .
Also using Eq. 3.88, we have {t y (t )} ( d ds ){ y (t )} ( d ds )Y ( s ) dY ds .
Example 3.151 [CH-2008 (2 marks)]: The Laplace transform of the function f (t ) t sin t is
(a) 2s ( s 2 1) 2 (b) 1 {s 2 ( s 2 1)} (c) {1 s 2 } {1 ( s 2 1)} (d) 1 {( s 1) 2 1}
Solution (a): From Eq. 3.88, {t sin t} ( 1)1 ( d ds){sin t} (d ds ) 1 ( s 2 12 ) 2s ( s 2 1) 2
Example 3.152 [EC-2012, EE-2012, IN-2012 (1 mark), EC-2014 (2 marks)]: The unilateral
Laplace transform of f (t ) is 1 ( s 2 s 1) . The unilateral Laplace transform of t f (t ) is
(a) s ( s 2 s 1) 2 (b) (2 s 1) ( s 2 s 1) 2 (c) s ( s 2 s 1) 2 (d) (2s 1) ( s 2 s 1) 2
Solution (d): { f (t )} 1 ( s 2 s 1) 2 , so using Eq. 3.88 with n 1 we have
2 2 2
{t f (t )} ( 1)(d ds ) F ( s) ( d ds){1 ( s s 1)} (2 s 1) ( s s 1) .
Integral of Transform: If { f (t )} F ( s ) then f (t ) t F ( s )ds , provided lim f (t ) t
s t 0
exists. Proof: We have F ( s ) { f (t )} e st f (t ) dt . Now integrating both sides w.r.t. ‘ s ’ by
0
0
e
f (t )dt ds . As s and t are independent
variables so by interchanging the order of integration we have,
s
F (s)
t 0
s
e st ds f (t ) dt {e st ( t )}s s f (t ) dt
t 0
t 0
e st { f (t ) t}dt f (t ) t .
Example 3.153 [CE-2007 (2 marks)]: Evaluate 0 (1 t ) sin tdt
(a) (b) 2 (c) 4 (d) 3
s 2
s 1
1
ds (tan 1 s) s
sin t
t
tan 1 tan 1 s
2
s e st
0
sin t
t
dt s . Now taking limit s 0 both sides we
2
sin t sin t sin t
get lim e st dt lim s lim e st dt lim s dt .
s 0 0 t s 0 2 0 s0 t s 0 2 0 t 2
0 1 (1 x ) dx 0 (sin x) x dx
2
Example 3.154 [CE-2016 (2 marks)]: The value of is
(a) 2 (b) (c) 3 2 (d) 1
0 1 (1 x ) dx (tan x)0 tan () tan (0)
2 1 1 1
Solution (b): 2 . The 2nd integral is
sin t 1 sin x
evaluated in Example 3.23, i.e. 0 dt . So 0 1 x 2 dx 0 dx .
t 2 x 2 2
sin 2 t
Example 3.155 [EE-2016 (2 marks)]: The value of the integral 2 dt is equal to
t
(a) 0 (b) 0.5 (c) 1 (d) 2
sin 2 t sin 2 t sin 2 t sin 2 t
Solution (d): I 2 dt 2 2 0 dt 4 0 dt , as is an even
t t
t t
4 sin k
function. Now let 2 t k 2 dt dk dt dk (2 ) , so I
0
dk . As shown in
k
sin k 4 sin k 4
Example 3.23, we have 0 k
dk
2
. So I
0 k
dk
2
2
Transform of Integrals: If functions f and g are piecewise continuous on the interval [0, ) ,
then a special product, denoted by f g , is defined by the integral
t
f g f ( ) g (t ) d (3.89)
0
and is called the convolution of f and g .
The convolution f g is a function of t. For e.g.,
t
et sin t e sin(t ) d (1 2)( sin t cos t et ) . Figure 3.8: Changing
0
order of integration from
The convolution of two function is commutative, i.e., f g g f or
to
t t
0 f ( ) g (t )d f (t ) g ( ) d .
0
It is not true that the integral of a product of functions is the product of the integrals. However, it
is true that the Laplace transform of the special product (Eq. 3.89) is the product of the Laplace
transform of f and g . This means that it is possible to find the Laplace transform of the
convolution of two functions without actually evaluating the integral. If f (t ) and g (t ) are
piecewise continuous on [0, ) and of exponential order, then
Convolution Theorem: { f g} { f (t )}{g (t )} F ( s )G ( s ) (3.90)
Proof: Let F ( s ) { f (t )} e s f ( ) d and G ( s ) { g (t )} e s f ( ) d
0 0
F (s )G (s ) 0
s
e f ( )d 0
s
e
f ( )d
0
s ( )
0 e f ( ) g ( )d d
F ( s )G ( s ) f ( ) d e s ( ) g ( ) d . Holding fixed, we let, t dt d
0 0
F ( s )G ( s ) f ( ) d e st g (t ) dt . In the t plane we are integrating over the shaded
0 0
region, as shown in Fig. 3.8. Since f and g are piecewise continuous on [0, ) and of
exponential order, it is possible to interchange the order of integration:
t
F ( s)G ( s ) e st dt f ( ) g (t ) d e st
0 0
0 f ( )g (t )d dt { f g}
t
Inverse form of Eq. 3.90: The convolution theorem is sometimes useful in finding the inverse
Laplace transform of the product of two Laplace transforms. Hence, from Eq. 3.90, we have
1{F ( s)G ( s )} f g (3.91)
Example 3.160 [EE-2007 (2 marks)]: If u (t ) , r (t ) denote the unit step and unit ramp functions
respectively and u (t ) r (t ) their convolution, then the function u (t 1) r (t 2) is given by
(a) (1 2)(t 1)(t 2) (b) (1 2)(t 1)(t 2) (c) (1 2)(t 1) 2 u (t 1) (d) None of these
Solution (c): Let g (t ) u (t 1) r (t 2) then from Eqs. 3.90 and 3.91, we have
{ g (t )} {u ( t 1) r (t 2)} {u (t 1)}{r ( t 2)} ( e s s ) ( e 2 s s 2 ) e s s 3
g (t ) 1{e s s 3 } . Now using second shifting theorem given by Eq. 3.92, i.e., we have
1{e a s F ( s )} f (t a) (t a ) where f (t ) 1{F ( s )} so 1 e s (1 s3 ) f (t 1) (t 1) ,
1 1
where f (t ) {1 s } (1 2!) {2! s } t 2 f (t 1) (1 2)(t 1) 2 .
3 3 2
Thus
1
e s 3
2
(1 s ) (1 2)(t 1) u (t 1) g (t )
Example 3.161 [EE-2011 (1 mark)]: Given two continuous time signals x(t ) e t and y (t ) e 2t
which exist for t 0 , the convolution z (t ) x(t ) y (t ) is
(a) e t e 2 t (b) e 3t (c) e t (d) e t e 2t
Solution (a): From Eqs. 3.90 and 3.91, we have {z (t )} {x (t ) y (t )} {x (t )}{ y (t )}
{z (t )} {e t }{e 2 t } {1 ( s 1)} {1 ( s 2)} {1 ( s 1)} {1 ( s 2)}
z (t ) 1 {1 ( s 1)} {1 ( s 2)} 1 1 ( s 1) 1 1 ( s 2) e t e 2 t .
0
t
Solution: With f (t ) et and g (t ) sin t , the convolution theorem states that the Laplace transform
of the convolution of f and g is the product of their Laplace transforms:
0
t 2 2
When g (t ) 1 {g (t )} {1} 1 s , the convolution theorem, i.e. Eq. 3.90, imples that the
Laplace transform of the integral of f is
f ( )d F (s) s
t
0
(3.92)
[Eq. 3.92 was asked in EC-2009, ME-2007 (2 marks)]; and the inverse of Eq. 3.92 is given as,
t 1
0 f ( )d {F ( s) s} (3.93)
which can be used in partial fraction when s n is a factor of the denominator and
f (t ) 1{F ( s )} is easy to integrate. For e.g. we know for f (t ) sin t that
f (t ) sin t F ( s) { f (t )} 1 ( s 2 1) and hence by Eq. 3.93,
t
1
1 {s( s 2
1)} 1
{1 (s 2
1)} s sin d 1 cos t
0
Example 3.163 [CE-2000 (2 marks)]: Let F ( s) [ f (t )] denote the Laplace transform of the
function f (t ) . Which of the following statement is correct?
1
(a) ( df dt ) (1 s ) F ( s ) ; f ( ) d sF ( s ) f (0)
0
1
(b) ( df dt ) sF ( s) F (0) ; f ( ) d dF ds
0
1
(c) ( df dt ) sF ( s) F (0) ; f ( ) d F ( s a)
0
1
(d) ( df dt ) sF ( s) F (0) ; f ( ) d (1 s ) F (s)
0
Solution (d): From Eq. 3.75 and Eq. 3.92 we can say that option (d) is correct.
Example 3.164 [CH-2014 (1 mark)]: For the time domain function f (t ) t 2 , which ONE of the
t
following is the Laplace transform of 0 f (t )dt ?
(a) 3 s 4 (b) 1 (4s 2 ) (c) 2 s 3 (d) 2 s 4
0
where
0
4
T st e s (T 2) e s (0) 1 e s (T 2)
1 e s (T 2)
e
0
u{t} u{t (T 2)} dt
s s s s s
1 T 1 1 e s (T 2)
1
Hence { f (t )} sT 0 e st u{t} u{t (T 2)} dt sT
.
1 e 1 e s s{1 e sT 2 }
Exercise: 3.3
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1 2 1 14 1 2 1 14
(c) F ( s) 2 3 2 e 6 s (d) F ( s) 2 3 2 e 6 s
s s s s s s s s
14. The solution of IVP y 10 y 9 y 5t , y (0) 1 , y (0) 2 is
(a) y (t ) (50 81) (5 9)t (31 81)e 9t 2e t (b) y (t ) (5 8) (5 9)t (3 8)e9t 2et
(c) y (t ) (50 81) (5 9)t (31 81)e 9t 2et (d) y (t ) (5 8) (5 9)t (3 8)e9t 2e t
15. The solution of IVP 2 y 3 y 2 y te 2 t , y (0) 0 , y (0) 2 is
(a) y (t ) (1 125){96et 2 96e 2t 10te 2t (25 2)t 2 e 2 t }
(b) y (t ) (1 125){96et 2 96e 2 t 10te 2 t (25 2)t 2 e 2 t }
(c) y (t ) (1 125){96et 2 96e 2 t 10te 2t (25 2)t 2 e 2t }
(d) y (t ) (1 125){96et 2 96e 2t 10te 2 t (25 2)t 2 e 2 t }
16. The solution of IVP ty ty y 2 , y (0) 2 , y (0) 4 is
(a) y (t ) 2 4t (b) y (t ) 2 4t (c) y (t ) 2 4t (d) y (t ) 2 4t
17. The solution of IVP y 2 y 15 y 6 (t 9) , y (0) 5 , y (0) 7 is
(a) y (t ) (3 4)u (t 9){e3(t 9) e 5( t 9) } (9 4)e3t (11 4)e 5 t
(b) y (t ) (3 4)u (t 9){e 3(t 9) e 5( t 9) } (9 4)e3t (11 4)e 5t
(c) y (t ) (3 4)u (t 9){e3(t 9) e 5( t 9) } (9 4)e 3t (11 4)e 5t
(d) y (t ) (3 4)u (t 9){e3(t 9) e 5( t 9) } (9 4)e3t (11 4)e 5t
18. The Laplace transform of the periodic function
(square wave) shown in figure is
(a) 1 {s (1 e s )} (b) 1 {s(1 e s )}
(c) 1 {s(1 e s )} (d) 1 (1 e s )
2
19. Given two continuous functions x(t ) 1 t and y (t ) t which exist for t 0 , the convolution
z (t ) x(t ) y (t ) at t 4 is _____.
x
20. The solution of the integral equation xe x e x t y (t ) dt is
0
x
(a) y ( x ) e x
(b) y ( x ) e (c) y ( x ) e1 x (d) y ( x) e 1 x
Fourier Analysis: Let f (t ) be a piecewise continuous function (i.e. having finite number of
removable or jump discontinuities) on [ T , T ] . Then the Fourier series of f (t ) is the series given as:
f (t ) a0 n 1 ( an cos nt bn sin n t ) (3.95)
th
where 2 (2T ) is the fundamental frequency; n is the n harmonic of fundamental
frequency; the coefficients a0 , an , bn are called Fourier coefficients of f (t ) , which are defined as:
T T T
a0 {1 (2T )} f (t ) dt , an (1 T ) f (t ) cos( nt ) dt , bn (1 T ) f (t ) sin( nt ) dt (3.96)
T T T
Proof: If we integrate both sides of the Eq. 3.95 and assume that it is permissible to integrate the
T T T
series term by termT f (t )dx T a0 dx T n 1{an cos(nt ) bn sin(nt )}dx
we get,
T T T T
f (t ) dt a0 dt n 1 an cos( nt ) dt n 1 bn sin( n t )dt a0 (2T )
T T T T
T T T
a0 {1 (2T )} f (t ) dt . As, T cos(nt )dt 0 and T sin(nt )dx 0 for all n . To determine
T
an for n 1 , multiply Eq. 3.95 by cos(mt ) (where m I and m 1 ) and integrate from T to T ,
T T T
i.e., I
T
f (t ) cos( m t ) dt a0 cos( mt ) dx
T T
n 1
( an cos( nt ) bn sin( nt ) cos( mt ) dt
T T T
I a0 cos( m t )dt n 1 an cos( nt ) cos( m t ) dt n 1 bn sin( nt ) cos( m t ) dt . As the
T T T
first integral in the RHS is shown to be zero and also using trigonometric identities we can prove that
T T 0, n m
T sin( n t ) cos( m t ) dt 0 for all n and m ; also T cos( n t ) cos( m t ) dt . So
T , n m
T T T
f (t ) cos( m t ) dt amT am (1 T ) f (t ) cos(mt ) dt an (1 T ) f (t ) cos(nt ) dt ,
T T T
where n 1, 2, 3, . Similarly, if we multiply Eq. 3.95 by sin(mt ) and integrating it from T to T ,
T
we get bn (1 T ) f (t ) sin( nt ) dt , where n 1, 2, 3, .
T
The Fourier coefficients of a sum f1 f2 are the sums of the corresponding Fourier coefficients
of f1 and f 2 .
The Fourier coefficients of cf are c times the corresponding Fourier coefficients of f .
Example 3.169
(a) f ( x ) (1 2) n 1{2 ( n )}sin( n 2)
[CE-2000 (1
mark)]: A function (b) f ( x ) n 1{2 (n )}sin(n 2) cos nx
with period 2 is
shown below. The (c) f ( x ) (1 2) n 1 {2 (n )}sin( n 2)
Fourier series for
(d) f ( x ) n 1 {2 ( n )}sin( n 2) sin nx
this function is given by
Solution (b): As f ( x ) is periodic with period 2T 2 T and from the given figure, we have
0, x 2
f ( x ) 1, 2 x 2 . So from Eq. 3.95, the Fourier series, with 2 (2T ) 1 , is given as
0, 2 x
f ( x ) a0 n 1 ( an cos nx bn sin nx ) , where from Eq. 3.96
a0 {1 (2 )}
f ( x )dx {1 (2 )}
2
(0)dx
2
2
(1)dx
2
(0) dx 1 2 ;
f ( x) cos nxdx (1 )
2 2
an (1 ) (0) cos nxdx 1cos nxdx (0) cos nxdx
2 2
2
an (1 ) 0 (sin nx) n 2 0 {1 (n )} 2 sin( n 2) {2 n }sin(n 2) ;
bn (1 )
f ( x ) sin nxdx (1 ) 2
(0) sin nxdx
2
2
1sin nxdx
2
(0) sin nxdx
bn (1 ) 0 (cos nx ) n 2 0 0 .
2
Thus Fourier series for the given figure is f ( x ) (1 2) n 1{2 ( n )}sin( n 2) cos nx
Example 3.170 [EC-2005 (1 mark)]: Choose the function f (t ) ; 1 , for which a Fourier
series cannot be defined.
(a) 3sin(25t ) (b) 4 cos(20t 3) 2 sin(10t ) (c) e t sin(25t ) (d) 1
Solution (c): Functions given in options (a), (b) and (d) are periodic so we can find their Fourier
series. But the function given in option (c) is not periodic so we cannot find its Fourier series.
2 0 2 12
a0 x (1 2) x 1 2
(1 2) x 0
0 (1 2) (1 2)(1 2) 2 (1 2)(1 2) 2 0 1 2
3.4.1 Waveform Symmetry
There are few methods by which the evaluation of Fourier coefficients is simplified by symmetry
consideration. These methods reduce the amount of calculations involved in finding out the
coefficients. From Eq. 3.96, we have a0 {1 (2T )}
T
T
f (t )dt {1 (2T )} 0
T
T
f (t ) dt f (t ) dt , now
0
putting t dt d in the first integrand on the RHS, we get
0 T
a0 {1 (2T )} f ( ) d f (t ) dt {1 (2T )}
T 0 T
0 0
T
f ( )d f (t ) dt
T
a0 {1 (2T )} { f (t ) f ( t )}dt . Again from Eq. 3.96, we have
0
an (1 T )
T
T
f (t ) cos(nt )dt (1 T ) 0
T
T
f (t ) cos(nt )dt f (t ) cos(nt )dt .
0 Now putting
t dt d in the first integrand on the RHS, we get
0
T
T
an (1 T ) f ( ) cos(n )d f (t ) cos(nt )dt
0
an (1 T ) T
0
T
f ( ) cos(n )d f (t ) cos(nt )dt
0
a (1 T ) { f (t) f (t)}cos(nt)dt
T T T
n f ( t ) cos(nt )dt f (t ) cos(nt )dt an (1 T )
0 0 0
have a0 0 , an 0 and bn (2 T )
0
T
f (t ) sin(nt )dt . Thus the Fourier series expansion of an
odd function contains only the sine terms, the cosine and the constant terms being zero. The
Fourier series of odd symmetric function also called as Fourier Sine series. We also have
f (t ) f (2T t ) for odd symmetry which comes if an 0 , where 2T is the period of f (t ) .
Thus the Fourier series of a real periodic function has only cosine terms if it is even and sine
terms if it is odd. [This point was asked in EC-1994, EC-2009 (1 mark)].
half of the wave is exactly opposite to the first half. A function with half wave symmetry does not
have to be even or odd, as this property requires only the shifted signal is opposite. So the value of the
Fourier constants is given as:
a0 {1 (2T )}
T
T
f (t )dt {1 (2T )} 0
T
T
f (t ) dt f (t ) dt , so putting t T dt d in the first
0
0
f ( T ) d
T
0
f (t )dt {1 (2T )} { f ( )}d
T
0 0
T
f (t ) dt
T
a0 {1 (2T )} f (t )dt f (t )dt 0
0
T
0 . Now, an (1 T )
T
T
f (t ) cos( nt ) dt ,
an (1 T ) 0
T
T
f (t ) cos( n t )dt f (t ) cos(nt ) dt , let t T dt d in the first integrand
0
0 T T
we get I1 f (t ) cos( n t )dt f ( T ) cos{n ( T )}d f ( ) cos( n nT ) d
T 0 0
T T
I1 f ( ) cos( n n ) d f ( ){cos( n ) cos n sin( n ) sin n }d
0 0
T T
I1 f ( ){cos( n ) cos n }d cos n f (t ) cos(nt ) dt . Thus we have
0 0
T T
an (1 T ) cos n f (t ) cos(nt )dt f (t ) cos(nt )dt {(1 cos n ) T }
0 0 0
T
f (t ) cos(nt )dt
0, n is even 0, n is even
an T . Similarly bn T .
(2 T ) f (t ) cos( nt ) dt ,
0
n is odd (2 T ) f (t ) sin(nt ) dt ,
0
n is odd
Thus the Fourier series expansion of a periodic function having half wave symmetry contains only
odd harmonics and the constant term being zero.
Example 3.173 [CE-2003 (2 marks)]: The Fourier series expansion of a symmetric and even
1 (2 x ), x 0
function, f ( x ) where, f ( x ) and will be
1 (2 x ), 0 x
(a) n1{4 ( 2 n 2 )}(1 cos n ) cos nx (b) n1{4 ( 2 n 2 )}(1 cos n ) cos nx
(c) n1{4 ( 2 n 2 )}(1 sin n ) cos nx (d) n1{4 ( 2 n 2 )}(1 sin n ) cos nx
Solution (b): The Fourier series expansion of a symmetric and even function f ( x ) have only cosine
term. As the given function have period of 2 , so with 2T 2 , f ( x ) a0 n 1 ( an cos nx ) ,
1 2x 1 1 x2 1 2
where a0 g ( x ) dx 1 dx x 0 ; and
0 0 0
2 2 2x 2 2x
an g ( x) cos( nx) dx 1 cos( nx) dx cos nx dx cos nx dx
0 0 0 0
2 sin nx 2x 2 2x
an 0 cos nx dx 0 0 cos nx dx
n 0
2 2 x sin nx 2 2 cos nx
an
n 0
0
2 sin nx
n
dx
0
n n 0
2 2 cos n 1
n n
an {4 ( n 2 2 )}(1 cos n ) . Thus f ( x ) n 1 {4 ( n 2 2 )}(1 cos n ) cos nx
Example 3.174 [EE-2006 (1 mark)]: x (t ) is a real valued function of a real variable with period T .
Its trigonometric Fourier Series expansion contains no terms of frequency 2 (2 k ) T ; k 1, 2,
. Also, no sine terms are present. Then x (t ) satisfies the equation
(a) x (t ) x (t T ) (b) x (t ) x (T t ) x ( t )
(c) x (t ) x (T t ) x (t T 2) (d) x(t ) x (t T ) x (t T 2)
Solution (c): Please note that the formulas are derived for period 2T but in given question we have
period T , so we put T 2 in place of T . As the Fourier expansion of the function x (t ) contains no
frequency term containing even n so it means that it contains only odd n , i.e., odd harmonics so we
can say that x (t ) is a half wave symmetric function which means x(t ) x(t T 2) . Also the Fourier
expansion of x (t ) contains no sine term so x (t ) has is an even symmetry and thus we have
x (t ) x (T t ) . Thus x(t ) x(T t ) x(t T 2) .
Example 3.175 [EE-2005 (2 marks)]: The Fourier series for the function f ( x ) sin 2 x is
(a) sin x sin 2 x (b) 1 cos 2x (c) sin 2 x cos 2 x (d) 0.5 0.5cos 2x
2 2
Solution (d): As sin x (1 2) (1 2) cos 2 x so Fourier series of sin x is (1 2) (1 2) cos 2x .
1, T 4 t 3T 4
Example 3.176 [EE-2007 (2 marks)]: A signal x (t ) is given by x (t ) 1, 3T 4 t 7 T 4
x (t T ) otherwise
. Which among the following gives the fundamental Fourier term of x (t ) ?
4 t 4 t 4 t t 4
(a) cos (b) cos (c) sin (d) sin
T 4 2T 4
T 4 2T 4
Solution: The given function x (t ) is periodic with period 2T . Also x(t ) x (t T ) so x(t ) has half
wave symmetry so we have a0 0 . Now the derived formulas are for functions from T to T . But
we have from T 4 to 7T 4 so we shift the given function left side by 3T 4 and thus
1, T t 0
x (t 3T 4) . As we have to find the fundamental Fourier term so we have n 1
1, 0t T
and also as t {2 (2T )}T ; thus the Fourier coefficient
T T T
an (2 T ) f (t ) cos(t ) dt (2 T ) cos( t )dt {2 (T )} (sin t ) 0
0 0
T T
an {2 (T )}(sin T 0) 0 . Now bn (2 T ) f (t ) sin(t ) dt (2 T ) sin( t ) dt
0 0
T
bn (2 T ) (cos t ) 0 {2 (T )}(cos T 1)T0 4 .
Thus we have x t (3T 4) (4 ) sin t x (t ) (4 ) sin t (3T 4)
x (t ) (4 ) sin ( t T ) (3 4) (4 ) sin ( 4) ( t T )
x (t ) (4 ) sin ( 4) ( t T ) (4 ) sin ( t T ) ( 4) .
Quarter wave Symmetry: If a function is both half wave symmetry and even or odd symmetry
then it is considered as quarter wave symmetry, as shown in Fig. 3.10.
Figure 3.10: Quarter Wave Symmetric (a) Even (b) Odd Function
So for the quarter wave odd symmetric function the Fourier coefficients are a0 0 , an 0 and
T 2
bn (4 T ) f (t ) sin( nt ) dt , where n is odd. Also for the quarter wave even symmetric function
0
T 2
the Fourier coefficients are a0 0 , bn 0 and an (4 T ) f (t ) cos(nt ) dt , where n is odd.
0
2
2
( x) cos(m x 2)dx ( x) cos(m x 2)dx
0
0 2
1 m x
2 2
2 m x 2 m x 2 m x 2
am x sin 2 2 cos x sin 2 2
cos
m
2 2 m 2 2 m 2 m 2 0
1 22 22 5 x
am 2 2 2 1 cos m 2 2 1 cos m . Thus coefficient of cos is the value
2 m m 2
of am at m 5 a5 {2 2 (52 2 )}( 1 cos 5 ) {2 2 (5 2 2 )}( 2) 8 (5 2 2 )
[Similar question was also asked in TF-2009 (1 mark)]
Common Data for Questions 3.178 & 3.179: Let f : R R be defined by f ( x) x 2 for
x and f ( x 2 ) f ( x) .
2 2 sin nx sin nx 2 cos nx cos nx
I x cos(nx ) dx x 0 2 x dx 0 x dx
0 0
n 0 n n n 0 n
2 cos n sin nx 2 cos n 4 cos n 4
I 2 . So an 2
( 1) n . Now
n n n 0 n2 n n2
2
bn (1 ) f ( x) sin(nx ) dx (1 ) x 2 sin( nx ) dx 0 , as x sin( nx ) is an odd function. So
2 n 2
(1) 4 n cos nx
f ( x) 2
cos nx 4 ( 1) 2 .
3 n 1 n 3 n 1 n
Example 3.179 [XE-2007 (2 marks)]: The sum of absolute values of the Fourier coefficients of f is
(a) 2 6 (b) 2 3 (c) 2 2 3 (d) 2
Solution (d): As in above problem we have a0 2 3 , an ( 1) n (4 n 2 ) and bn 0 . So sum of
2
4
absolute values of the Fourier coefficients of f is S a0 an bn 0
3 n 1 n2
2
1 1 1 2 2 1 1 1 2
S 4 2 2 2 4 2 , as 2 2 2 .
3 1 2 3 3 6 1 2 3 6
Example 3.180 [XE-2013 (1 mark)]: The Fourier series of the periodic function f ( x) x ,
1 x 1 , f ( x 2) f ( x) , x R is given by (1 2) n 1{4 cos(2n 1) x} {(2 n 1) 2 2 } . Using
the above, the sum of the infinite series 1 (1 32 ) (1 52 ) is
(a) 2 4 (b) 3 2 8 (c) 2 8 (d) 2 2
Solution (c): From the given data, the given Fourier series can be written as
1 4 cos(2n 1) x 1 4 4 4
f ( x) 2 2
2 2 cos x 2 2 cos 3 x 2 2 cos 5 x . Now at
2 n 1 (2n 1) 2 1 3 5
1 4 1 1 1
x 0 we have f ( x ) 2 2
2 2 . As the given function f ( x ) , which is of period 2
2 1 3 5
x, 1 x 0 1, 1 x 0
, can be written as: f ( x ) f ( x ) f ( x ) and f ( x ) are
x, 0 x 1 1, 0 x 1
piecewise continuous in [ 1,1] . So the value f ( x ) at x 0 is the average of LHL ( f (0 ) ) and RHL
of f ( x) at x0 as x0 is the point of discontinuity. Thus
f (0 ) lim f (0 h) lim{(0 h)} 0 and f (0 ) lim f (0 h) lim(0 h) 0 . So
h 0 h0 h 0 h 0
2
f (0 ) f (0 ) 4 1 1 1 1 0 1 1 1
f ( x)
0 . Thus
.
2 2 2 12 32 52 12 32 52 8
[Similar questions were also asked in CS-1993, EC-1993 (1 mark)]
2, x 0
Example 3.181 [AE-2014 (2 marks)]: For the periodic function given by f ( x )
2, 0 x
with f ( x 2 ) f ( x) using Fourier series, the sum s 1 (1 3) (1 5) (1 7) converges to
(a) 1 (b) 3 (c) 4 (d) 5
Solution (c): As the given function is periodic with period 2 . So from Eqs. 3.95 and 3.96 with
2T 2 , we have f ( x ) a0 n 1 an cos nx bn sin nx , where
a0 {1 (2 )}
(2)dx (2)dx {1 (2 )}(2x)
f ( x)dx {1 (2 )}
0
0
0
(2 x)0 0 ,
0
an (1 ) 2 (sin nx ) n 2 (sin nx ) n 0 0 , similarly
bn (1 )
f ( x)sin(nx)dx (1 )
0
( 2) sin(nx)dx (2) sin(nx)dx
0
0
2 cos nx cos nx 2 1 cos n 1 cos n 4 1 cos n
bn
n n 0 n n n
4 1 cos n 42 2 2
Thus f ( x ) 0 n 1 sin nx sin x sin 3 x sin 5 x . So at x 2 ,
n 1 3 5
2 2 42 8 1 1 1
we have f ( x ) . Now as f ( x) and f ( x ) both are piece-
1 3 5 1 3 5
wise continuous function in [ , ] so the Fourier series converges whose value is f ( 2) , as
x 2 is a continuous point for f ( x) . Thus f ( 2) 2
8 1 1 1 1 1 1 2
2 . So the given series converges to .
1 3 5 1 3 5 8 4 4
0, x 0
Example 3.183 [CE-2016 (1 mark)]: The Fourier series of the function, f ( x )
x, 0 x
2 cos x cos 3 x sin x sin 2 x sin 3x
in the interval [ , ] is f ( x ) 2 2
. The
4 1 3 1 2 3
convergence of the above Fourier series at x 0 gives
1 2
( 1) n 1 2
1 2
(1) n 1
(a) 2 (b) (c) (d)
n 1 n 6 n 1 n2 12 n 1 (2 n 1)
2
8 n 1 (2 n 1) 4
2 cos 0 cos 0 sin 0 sin 0 sin 0
Solution (c): At x 0 , f (0) 12
2
4 3 1 2 3
2 2
21 1 1 1 1
0 2
2 0 2 2 n 1 .
4 1 3 1 3 8 (2n 1) 2 8
T
coefficients cn and kn are combined by writing kn c n then cn {1 (2T )} f (t )e int dt ,
T
n 0, 1, 2, . So the complex form of Fourier series given by Eq. 3.95 is given by Eq. 3.99.
Example 3.184 [EE-2008 (2 marks)]: Let x (t ) be a periodic signal with time period T . Let
y (t ) x (t t0 ) x (t t 0 ) for some t0 . The Fourier series coefficients of y (t ) are denoted by b . If
bk 0 for all odd k , then t0 can be equal to
(a) T 8 (b) T 4 (c) T 2 (d) 2T
Solution (b): Let ck be the coefficient of Fourier series of x (t ) then From Eq. 3.99,we have
ik ( t t0 ) ik ( t t0 )
x (t ) k ck e ikt x (t t 0 ) k ck e also x (t t0 ) k ck e . So
ik ( t t0 ) ik ( t t0 ) ik ( t t0 ) ik ( t t0 )
y (t ) x (t t 0 ) x (t t0 ) k ck e k ck e k ck {e e }
ik t0 ik t0 ikt0 ik t0
y (t ) k ck {eikt e eikt e } k ck (e e )eikt . So coefficient of Fourier
ik t0 ik t0
series of y (t ) is given as bk ck {e e } 2 ck cos( k t0 ) . As time period of x (t ) is T so
2 T . Now bk 0 for all odd k and so cos(k t0 ) 0 kt0 k ( 2) t0 { (2 )} T 4
Proof: Consider the complex form of Fourier series f (t ) n cn eint , where
cn {1 (2T )}
T
T
f (t )e int dt , n 0, 1, 2, . So, f (t ) n {1 (2T )}
T
T
f (t )e int dt eint . If
the period 2T becomes infinite, the function does not repeat itself and becomes non-periodic. So the
interval between the adjacent harmonic frequencies ( n 1) n (2 ) (2T ) , so we have
1 (2T ) (2 ) (2 ) . As T , d and the frequency goes from a discrete variable
over to a continuous variable. So, as T , 1 (2T ) d (2 ) and n ; and also the
summation sign change to integration and thus we have f (t ) {1 (2 )}
e
i t
f (t )e it dt d .
So the expression F (i ) f (t )e it dt is called the Fourier transform of f (t ) and the expression
f (t ) {1 (2 )} F (i )eit d is called the inverse Fourier transform.
it0
Time Shifting: If F { f (t )} F ( ) and t0 , then F { f (t t 0 )} e F ( )
Proof:
it0 i t0 i t0
F { f (t t0 )} f (t t0 )e i t dt e f (u )e
iu
du e f (t )e
it
dt e F ( )
i0
Frequency Shifting: If F { f (t )} F ( ) and , then F ( 0 ) F {e f (t )}
i0 i0 i ( 0 ) t
Proof: F {e f (t )} e f (t )e i t dt f (t )e dt F ( 0 )
Symmetry: If F { f (t )} F ( ) then F {F (t )} 2 f ( )
Proof:
f (t ) F 1{F ( )} {1 (2 )} F ( )e it d {1 (2 )} F ( )eit d {1 (2 )} F ( x )eixt dx
and f ( ) {1 (2 )} F ( x )e ix dx 2 f ( ) F (t )e it dt F{ f (t )}
Differentiation in Time: Let n N and suppose that f ( n ) is piecewise continuous function and
lim f ( k ) (t ) 0 then F { f ( n ) (t )} (i ) n F ( ) .
t
e t , t 0
Example 3.185 [IN-2011 (1 mark)]: Consider the signal x (t ) . Let X ( ) denote the
0, t 0
Fourier transform of this signal. The integral {1 (2 )} X ( ) d is
F (i ) { A ( a i )} { A ( a i )} .
a t a t
So Fourier transform of f (t ) Ae is F { Ae } F (i ) { A ( a i )} { A ( a i )} .
at
f (t ) Ae u (t ) , where u (t ) is the unit step function. Note that Fourier transform will exists if
a0.
F (i ) F { f (t )} f (t )e i t dt A e at u (t )e it dt A e at e it dt
0
F (i ) A e ( a i ) t dt { A ( a i )}{e ( a i ) t }0 { A ( a i )}{e ( a i ) t }0 A ( a i ) .
0
a t i t 2ka 0, 0
F (i ) F { f (t )} lim ke e dt lim 2 2
, which is an impulse
a 0 a 0 a , 0
function, whose magnitude is found by F{k}d {2ka ( a 2 2 )}d 2 k .
So the Fourier transform f (t ) k is an impulse function at 0 whose magnitude is 2 k and
so F {k} F (i ) 2 k ( ) .
f (t ) (t ) , where (t ) is a unit impulse function or dirac delta function.
We know that the product of any arbitrary function f (t ) with unit impulse function (t t0 )
provides the function (t t0 ) to exists only at t t0 . So f (t ) (t t0 ) f (t 0 ) . Now
F { (t )} F (i ) (t )e it dt e i (0) 1 . So Fourier transform of unit impulse function is 1.
Relationship between Fourier Transform and Laplace Transform: For a function f (t ) , the
Laplace transform is defined as F ( s) f (t )e st dt and Fourier transform is defined as
0
it
F (i ) f (t )e dt . Laplace transform is one sided in the interval 0 t . Thus Laplace
transform is applicable for positive time function, f (t ) , t 0 ; while Fourier transform is applicable
for functions defined for all times. Laplace transform includes the initial conditions and is applicable
for transient analysis; while Fourier transform is only applicable for steady-state analysis. We know
that the Fourier transform of a function f (t ) is given as F (i ) f (t )e it dt ; also the Fourier
transform can be calculated only if f (t ) is absolutely integrable, i.e., f (t ) dt . If s i ,
so F ( s) f (t )e ( i ) t dt { f (t )e t }e i t dt Laplace transform of f (t ) is basically the
t
Fourier transform of f (t )e . If 0 then s i , then we have F ( s) f (t )e i t dt
F ( s ) F (i ) when s i . This means Laplace transform is same as Fourier transform when
s i . So Laplace transform provides broader characterization compared to Fourier transform.
s i indicates imaginary axis in complex s plane.
0, t 0
Solution (b): F (i ) F{x (t )} x (t )e i t dt e at u ( t )e it dt . As u (t )
1, t 0
0, t 0 0
u (t ) . F (i ) e ( a i ) t dt {1 ( a i )}{e ( a i ) t }0
1, t 0
F (i ) {1 ( a i )}{e 0 lim e ( ai )t } . Now lim e ( ai )t exists if a i 0 or a 0 .
t t
j t 2
Example 3.188 [EC-2014 (2 marks)]: For a function g (t ) , it is given that g (t )e dt e 2
t
for any real value . If y (t ) g ( ) d , then y (t )dt is
0, t 0, t 0, t
Now as u{( t )} u (t ) g ( )u{( t )}
1, t 1, t g ( ), t
0, t
g ( )u{( t )} .
g ( ), t
t
Thus we have y (t ) g ( ) d g ( )u{( t )}d g ( )u (t ) d which is the
convolution of g (t ) and u (t ) . So F { y (t )} F {g (t ) u (t )} F {g (t )} F {u (t )}
2
2 1 2 1 2 2 e 2
F { y (t )} e 2 (t ) y (t )e i t dt e 2 (t ) e (t ) .
i i i
i (0) t 2 2
Thus at 0 we have y(t )e dt (0)e 2(0) (t ) (1 i )e 2(0) 1 i i .
Exercise: 3.4
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1 1 x 0
1. The Fourier series for f ( x ) 1 2 x 0 for x [ 1,1] is
x 0 x 1
3 ( 1) n 1 1
(a) f ( x ) 2 2 cos(n x) sin( n x )
4 n 1 n n
3
( 1) n 1
(b) f ( x ) 2 2
cos( n x ) sin(n x)
4 n n 1 n
n
3 ( 1) 1
(c) f ( x ) 2 2 cos( n x ) sin(n x)
4 n 1 n
n
3 ( 1) 1 1
(d) f ( x ) 2 2 cos( n x ) sin( n x)
4 n 1 n n
2. Let f ( x ) be the function on [ 3, 3] which is graphed in the
figure. The constant term in the Fourier series for f is
(a) 5
(b) 5 2
(c) 3
(d) 3 2
0, x 2
3. Let f ( x ) 1, 2 x 2 , the Fourier series for f on the interval [ , ] is
0, 2 x
1 cos x cos(3x ) cos(5 x ) 1 2 cos x 2 cos(3 x ) 2 cos(5 x )
(a) (b)
2 3 5 2 3 5
1 cos x cos(3 x ) cos(5 x ) 1 2 cos x 2 cos(3 x ) 2 cos(5 x )
(c) (d)
2 3 5 2 3 5
4. Let f ( x ) be the function of period 2T 2 which is given on the interval ( 1,1) by f ( x ) 1 x 2
, the Fourier series of f ( x ) is
2 4
cos( n x ) 2
( 1) n 1
(a) f ( x )
3
2
n 1 n2
(b) f ( x )
3
n2
cos( n x)
n 1
2 4
(1) n 1 2 4
(1) n
(c) f ( x )
3 2 n 1 n 2
cos(n x) cos(n x)
(d) f ( x )
3 2 n 1 n 2
5. For the function f ( x ) 1 x defined on the half period 0 x 2 , the even extension to the full
period 2 x 2 is
(a) f ( x ) 1 x , 2 x 0 (b) f ( x ) 1 x , 2 x 0
(c) f ( x ) x , 2 x 0 (d) f ( x ) x 1 , 2 x 0
x 5, 0 x 1
6. For the function f ( x ) , the even extension to the full period 2 x 2 is
0, 1 x 2
0, 2 x 1 x 5, 2 x 1
(a) f ( x ) (b) f ( x )
x 5, 1 x 0 0, 1 x 0
0, 2 x 1 0, 2 x 1
(c) f ( x ) (d) f ( x )
x 5, 1 x 0 x 5, 1 x 0
7. For the function f ( x ) sin x defined on the half period 0 x , the odd extension to the full
period x is
(a) f ( x ) sin x , x 0 (b) f ( x ) sin x , x 0
(c) f ( x ) cos x , x 0 (d) f ( x ) cos x , x 0
0, 0 x 1
8. For the function f ( x ) , the odd extension to the full period 2 x 2 is
1 x, 1 x 2
1 x, 2 x 1 0, 2 x 1
(a) f ( x ) (b) f ( x )
0, 1 x 0 1 x , 1 x 0
x 1, 2 x 1 1 x , 2 x 1
(c) f ( x ) (d) f ( x )
1 x 0
0, 0, 1 x 0
9. Two dimensional Fourier transform and its inverse are infinitely
(a) aperiodic (b) periodic (c) linear (d) non-linear
at
e , t 0
10. The Fourier transform of the function f (t ) for a 0 is
0, t0
1 1 1 1
(a) (b) (c) 2 2
(d)
a i a i a i a i 2 2
u u v v v v u u
dividing (ii) by (iii) we get p q p q
x z y z x z y z
(u x u z p )(v y vz q ) (v x vz p )(u y u z q ) 0 , where u x u x , u y u y , …. Thus
(u , v ) ( u , v) (u , v )
( vz u y u z v y ) p (v x u z u x v z )q (u x v y v x u y ) p q Pp Qq R ,
( y, z ) ( z , x) ( x, y)
where (u , v ) ( y , z ) (vz u y u z v y ) P , (u , v ) ( z , x ) (vx u z u x vz ) Q and
(u , v ) ( x, y ) (u x v y vx u y ) R which are functions of x , y and z and do not contain p and
q . The equation Pp Qq R is called Lagrange’s linear equation. In fact, it is a quasi-linear
equation as the dependent variable may be present in P , Q and R . If P and Q are independent
of z and R is a linear in z the equation Pp Qq R is called linear.
Example 3.193 [XE-2016 (1 mark)]: Which of the following is a quasi-linear partial differential
equation?
2 2 2 4 3
2u 2 u u u u u u
(a) u 0 (b) 0 (c) 0 (d) 0
t 2 t x t x t x
Solution (a): A PDE is called a quasi-linear if all the terms with highest order derivatives of
dependent variables occur linearly, i.e. the coefficients of such terms are functions of only lower order
derivatives of the dependent variables. So among given PDEs the PDE ( 2u t 2 ) u 2 0 is quasi-
linear PDE and others are non-linear. So option (a) is correct.
PDE solvable by direct integration: Let us consider PDE which can be solved by direct
integration. While carrying out integration w.r.t. a variable the other variable is held fixed. So, in
place of constant of integration we have to add an arbitrary function of the variable held fixed.
Solution to Linear, Semi-Linear and Quasi-Linear PDE: If the PDE is linear or semi-linear
or quasi-linear then we can solve them by the method of Lagrange’s linear equation which is
described as follows: The equation of the form Pp Qq R …(i) is known as Lagrange’s equation,
where P , Q and R are functions of x , y and z . To solve this equation it is enough to solve the
subsidiary equations dx P dy Q dz R …(ii). If the solution of the subsidiary equation is of the
form u ( x, y ) c1 and v ( y , z ) c2 , then the solution of the given Lagrange’s equation is f (u, v) 0 .
Method of grouping: For solving the subsidiary equation dx P dy Q dz R ; take any two
members say first two or last two or first and last members. Now consider the first two members
dx P dy Q . If P and Q contain z (other than x and y ) try to eliminate it. Now direct integration
u ( x, y ) c1
gives Similarly take another two members dy Q dz R . If Q and R contain x (other
.
v ( y , z ) c2
than y and z ) try to eliminate it. Now direct integration gives . Therefore solution of the
given Lagrange’s equation is f ( u , v ) 0 . We will discuss various cases which are as follows:
Second Order Partial Differential Equation: As for ODEs, 2nd order PDEs will be the most
important ones in application. If we let u denote the dependent variable and let x and y denote the
independent variables, then the general form of a linear second – order PDE is given by
2u 2u 2u u u
A 2 B C 2 D E Fu G (3.105)
x xy y x y
where, the coefficients A, B , , G are functions of x and y ; and at least one of the coefficients A ,
B and C is not zero. When G ( x, y ) 0 , Eq. 3.105 is said to be linear and homogeneous; if G 0
then Eq. 3.105 is said to be linear and non – homogeneous. [This point was asked in AG-2014 (1
2 u 2u 2u u
mark)]. For e.g., the linear equations 0 and xy are homogeneous and non –
x 2 y 2 x 2 y
homogeneous equations, respectively.
Classification of Equation: A linear second-order PDE (Eq. 3.105) in two independent variables
with constant coefficients, such that at least one of the coefficients A , B and C is not-zero, is:
Hyperbolic if B 2 4 AC 0 Parabolic if B 2 4 AC 0 Elliptic if B 2 4 AC 0
Example 3.198 [AE-2010 (1 mark)]: The linear second order partial differential equation
2 2 2
5 2 3 2 2 9 0 is
x xy y
(a) Parabolic (b) Hyperbolic (c) Elliptic (d) None of the above
Solution: Comparing the given equation with Eq. 3.107, we have A 5 , B 3 , C 2 ; as
2
B 4 AC 9 40 31 0 the given PDE is of elliptic type.
u 2u
Example 3.199 [AE-2016 (1 mark)]: The partial differential equation , where is a
t x 2
positive constant, is
(a) circular (b) elliptic (c) hyperbolic (d) parabolic
2 2 2
u u u u u
Solution (d): Comparing the given PDE with A 2 B C 2 D E Fu G , we
x xt t x t
have B C D F G 0 and A , E 1 , so B 2 4 AC (0) 2 4( )(0) 0 , thus given
PDE is parabolic in nature.
[Similar questions were also asked in IN-2013, ME-1996 (1 mark)]
Example 3.200 [MN-2010 (1 mark)]: The partial differential equation, r ( r ) k , where k is
constant, is a solution for
2 1 2 2
2 1 2 1
(a) 0 (b) 0 (c) r 0 (d) 0
r 2 r r r 2 r r 2 r r r 2 r r
Solution (d): Partial differentiating the given solution w.r.t. r, we get
2 2
1
r (k ) r 2 0 0.
r r r r r r r r 2
Solution of a PDE: A solution of a linear PDE (Eq. 3.105) is a function u ( x, y ) of two independent
variables that possesses all partial derivatives occurring in the equation and satisfies the equation in
some region of the xy plane. Finding the general solutions of linear PDEs is often difficult and also
a general solution is not useful in applications; so our focus will be on finding particular solutions of
some important linear PDEs, i.e., equations that appear in many applications.
Superposition principle: If u1 , u2 , , uk are the solutions of a homogeneous linear PDE, then the
linear combination, u c1u1 c2 u2 ck uk , where c1 , c2 , , ck are constant, is also a solution.
Separation of variable method: In this method we seek a particular solution of the form of a
product of a function of x and a function of y , i.e., u ( x, y ) X ( x)Y ( y ) . With this assumption it
is sometimes possible to reduce a linear PDE in two variables to two ODEs.
Important second order PDEs: The following equation are some common PDEs:
2u 2 2u 2u 2u
1 D Wave equation: c 2 D Poisson equation: f ( x, y )
t 2 x 2 x 2 y 2
u 2 2u 2u 2 2u 2u
1 D Heat equation: c 2 D Wave equation: 2
c 2
2
t x 2 t x y
2u 2u 2u 2u 2u
2 D Laplace equation: 0 3 D Laplace equation: 0
x 2 y 2 x 2 y 2 z 2
We are going to discuss: One dimensional Wave equation (which is Hyperbolic); One dimensional
Heat equation (which is Parabolic); and Two dimensional Laplace equation (which is Elliptic).
Under these assumptions we may expect solutions u ( x, t ) that describe the physical reality
sufficiently well. To obtain the PDE, we consider the forces acting on a small portion of the string, as
shown in Fig. 3.12(b).
Since the string offers no resistance to bending, the tension is tangential to the curve of the string at
each point. Let T1 and T2 be the tension at the endpoints P and Q of that portion. Since the points of
the string move vertically, there is no motion in the horizontal direction. Hence the horizontal
components of the tension must be constant. So,
T1 cos 1 T2 cos 2 T (say) constant (3.106)
In the vertical direction the force balance T2 sin 2 T1 sin 1 x ( 2 u t 2 ) (3.107)
2 2
where the mass x of the portion times the acceleration u t , evaluated at some point between
x and x x ; here is the mass of the un-deflected string per unit length, and x is the length of
the portion of the un-deflected string. Now dividing Eq. 3.107 by Eq. 3.106
T sin 2 T1 sin 1 2u
2 tan 2 tan 1 x 2 (3.108)
T2 cos 2 T1 cos 1 t
where, tan 1 and tan 2 are the slopes of the string at x and x x , i.e., tan 1 (u x ) x and
tan 2 (u x ) x x . Here we write partial derivatives because u depends also on time t . Dividing
2
1 u u u . If we let x 0 , the linear PDE is,
Eq. 3.108 by x , we have 2
x x x x x x T t
2u 2 2u 2 T
2
c 2
, c (3.109)
t x
Eq. 3.109 is called the one – dimensional wave equation, which is homogeneous and of 2nd order.
The physical constant T is denoted by c 2 (instead of c ) to indicate that this constant is positive, a
fact that will be essential to the form of the solutions. Since the string is fastened at the ends x 0
and x L , we have the two boundary conditions:
u (0, t ) 0 and u ( L, t ) 0 for all t (3.110)
Also, form of the motion of the string will depend on its initial deflection (deflection at time t 0 ) as
f ( x ) ; and on its initial velocity (velocity at t 0 ), as g ( x ) . Hence we have two initial conditions:
u ( x, 0) f ( x) and (u t ) ( x,0) g ( x ) for all 0 x L , (3.111)
We now have to find the solution of Eq. 3.109 satisfying the conditions in Eq. 3.110 and 3.111. In the
method of separating variables, we determine solutions of the wave equation 3.109 of the form,
u ( x, t ) F ( x)G (t ) (3.112)
which are a product of two functions, each depending only one of the variables x and t .
2u 2G 2u 2F
Differentiating Eq. 3.112 we get F 2 and G 2 ; inserting these two in Eq. 3.109
t 2 t x 2 x
2 2 2 2
G 2 F 1 G 1 F
we get F 2 c G 2 2 k (say); here the separation constant k is still
t x c G t 2 F x 2
arbitrary. Now we have two single variable equations, i.e.
( 2 F x 2 ) kF 0 (3.113) ( 2G t 2 ) c 2 kG 0 (3.114)
We now determine solutions F and G of Eq. 3.113 and 3.114 so that u FG satisfies boundary
conditions (Eq. 3.110) , i.e.,
u (0, t ) 0 F (0)G (t ) 0 and u ( L, t ) 0 F ( L)G (t ) 0 for all t (3.115)
Now from Eq. 3.115, if G 0 , then u FG 0 , which is of no interest; hence G 0 and
F (0) 0 and F ( L) 0 (3.116)
For k 0 , the general of Eq. 3.113 is F ax b , and from Eq. 3.116 we have a b 0 , so that
F 0 u FG 0 , which is of no interest. For positive k 2 a general solution of Eq. 3.113 is
x x
F Ae Be and from Eq. 3.116, we get F 0 . Hence we are left with the possibility of
choosing k negative, say, k p 2 . Then Eq. 3.112 becomes ( 2 F x 2 ) p 2 F 0 and has a general
solution F ( x ) A cos px B sin px ; and after applying Eq. 3.116 we have, F (0) 0 A 0 and
F ( L) 0 B sin pL 0 sin pL 0 (since, if B 0 then F 0 ) pL n p n L ,
n 1, 2, ; for negative integer we get the same solution with minus sign.
Hence we must choose the Bn so that u ( x, 0) becomes the Fourier Sine Series of f ( x ) . Thus
L
Bn (2 L ) f ( x ) sin ( n L ) x dx , n 1, 2, (3.123)
0
Now by differentiating Eq. 3.121 w.r.t. t and using second initial condition of Eq. 3.111, we get
(u t ) ( x ,0) n1 Bn*n sin ( n L) x g ( x ) (3.124)
Hence we must choose the Bn* so that for t 0 the derivative u t becomes the Fourier Sine Series
of g ( x ) . Thus from Eq. 3.124, we have
L
Bn* {2 (cn )} g ( x ) sin (n L ) x dx , n 1, 2, (3.125)
0
Hence the solution of Eq. 3.109 with conditions given by Eq. 3.110 and 3.111 is given by Eq. 3.121
whose coefficients Bn and Bn* given by Eq. 3.123 and 3.124.
Example 3.204 [MA-2016 (1 mark)]: Let u ( x, t ) be the d’Alembert’s solution of the initial value
problem for the wave equation: utt c 2 u xx 0 , u ( x, 0) f ( x ) , ut ( x, 0) g ( x ) , where c is a positive
real number and f , g are smooth odd functions. Then u (0,1) is equal to _____.
Solution: The given problem is Cauchy Problem for the infinite string, D’Alembert’s solution. On
referring to Section 3.5.1, the solution of utt c 2 u xx 0 , u ( x, 0) f ( x ) , ut ( x, 0) g ( x ) is given by
x ct
u ( x, t ) (1 2) f ( x ct ) f ( x ct ) {1 (2c)} g ( s )ds .
x ct
c
So u (0,1) (1 2) f (0 c) f (0 c) {1 (2c)} g ( s) ds (1 2) f (c) f (c) {1 (2c)}(0) 0 ,
c
c
since g ( x ) is an odd function so c g (s)ds 0 ; also f ( x ) is an odd function so ‘ f ( x ) f ( x) ’.
Cauchy Problem for the infinite string, D’Alembert’s solution: When we have an infinite
string with no boundary, then we have the following Cauchy problem, i.e.
2u 2
2 u
c , t 0 , x with ICs u ( x, 0) f ( x) and (u t ) ( x,0) g ( x ) (3.126)
t 2 x 2
The solution of the wave equation for the infinite string is given as:
x ct
u ( x, t ) (1 2) f ( x ct ) f ( x ct ) {1 (2c )} g ( s )ds (3.127)
x ct
Example 3.205 [XE-2007 (1 mark)]: Let u ( x, t ) be the solution of the initial value problem
2u 2u
9 , t 0 , x , u ( x, 0) x 5 , (u t ) ( x ,0) 0 . Then u (2, 2) is
t 2 x 2
(a) 7 (b) 13 (c) 14 (d) 26
Solution (a): Comparing the given equation with Eq. 3.126, we have c 3 , and from given
conditions, we have f ( x ) x 5 and g ( x ) 0 . Thus the solution of the given PDE is from Eq.
x ct
3.127, i.e., u ( x, t ) (1 2) f ( x 3t ) f ( x 3t ) {1 (2c )} (0) ds
x ct
K u x ( x x, t ) u x ( x, t )
ut (3.133)
x
Now by taking limit of Eq. 3.133, we get {K ( )}u xx ut or by letting c 2 K ( ) thermal
diffusivity. Hence one dimensional heat equation is given as,
2
u 2 u 2 K
c 2
or u t c 2 u xx , c thermal diffusivity (3.134)
t x
For solving Eq. 3.134 we need boundary and initial conditions. Since the ends x 0 and x L of the
bar are kept at temperature zero, so the boundary conditions are:
u (0, t ) 0 , u ( L, t ) 0 for all t (3.135)
Also, the initial temperature in the bar at time t 0 is f ( x ) , so we have the initial condition
u ( x, 0) f ( x) , [ f ( x ) is given] (3.136)
Here we must have f (0) 0 and f ( L) 0 because of Eqs. 3.135. We shall determine a solution of
Eq. 3.134 satisfying Eqs. 3.135 and 3.136; one initial condition will be enough, as opposed to two
initial conditions for the wave equation. Substitution of a product u ( x, t ) F ( x)G (t ) into Eq. 3.134
gives F (G t ) c 2 ( 2 F x 2 )G . To separate the variables, we divide by c 2 FG , obtaining
1 G 1 2 F
k (3.137)
c 2G t F x 2
For k 0 and k 0 the only solution u FG satisfying Eqs. 3.135 is u 0 . For k p 2 we have
{1 (c 2 G )}(G t ) (1 F )( 2 F x 2 ) p 2 . Now we have two single variable equations, i.e.
2F 2
G
c 2 p 2G 0 (3.139)
2
p F 0 (3.138)
x t
The general solution of Eq. 3.138 is given as
F ( x ) A cos px B sin px (3.140)
From the boundary conditions (Eqs. 3.135) it follows that u (0, t ) F (0)G (t ) 0 and
u ( L, t ) F ( L)G (t ) 0 . Since G 0 u 0 , we require F (0) 0 , F ( L) 0 and get F (0) A 0
by Eq. 3.142 and then F ( L) B sin pL 0 with B 0 ; thus sin pL 0 p n L , n 1, 2, .
Setting B 1 , we thus obtain the following solutions of Eq. 3.138 satisfying Eqs. 3.135 as,
Fn ( x) sin(n x L ) , n 1, 2, (3.141)
For p n L , Eq. 3.139 (G t ) n2 G 0 , where n cn L . It has the general solution as
2
Gn (t ) Bn e n t , n 1, 2, , where Bn is a constant (3.142)
n x 2
Hence the functions un ( x, t ) Fn ( x )Gn (t ) Bn sin e n t , n 1, 2,
(3.143)
L
are the solution of the heat equation given by Eq. 3.134 satisfying Eqs. 3.135. These are Eigen-
functions of the problem, corresponding to Eigen-values n cn L .
To obtain a solution that also satisfies the initial condition Eq. 3.136, we consider a series of these
Eigen-functions,
n x n2t cn
u ( x, t ) n 1 u n ( x, t ) n 1 Bn sin e , n (3.144)
L L
From this and Eq. 3.136, we have u ( x, 0) n 1 Bn sin( n x L ) f ( x ) . Hence for Eq. 3.144 satisfy
Eq. 3.136, the Bn ’s must be the coefficients of the Fourier Sine Series, as given by
L
Bn (2 L ) f ( x ) sin{( n L ) x}dx , n 1, 2, (3.145)
0
Hence the solution of Eq. 3.134 with conditions given by Eq. 3.135 and 3.136 is given by Eq. 3.144
whose coefficients Bn is given by Eq. 3.145.
Example 3.208 [AE-2008 (2 marks)]: Suppose the non-constant functions F ( x ) & G (t ) satisfy
( d 2 F dx 2 ) p 2 F 0 , ( dG dt ) c 2 p 2G 0 , where p and c are constants. Then the function
u ( x, t ) F ( x)G (t ) definitely satisfies
2u 2 2u u 2 2u 2u 2 2
(a) c (b) c (c) 2 u 0
c u 0 (d)
t 2 x 2 t x 2 t 2
Solution (b): As the given equations are Eqs 3.138 and 3.139 which are generated from Eq. 3.134,
whose general solution is given by Eq. 3.112. So option (b) is the correct answer.
(time independent) problems u t 0 , heat equation becomes Laplace’s equation which is given as,
2u 2u
0 2u (3.146)
x 2 y 2
A heat problem then consists of this PDE to be considered in some region R of the xy plane and a
given boundary condition on the boundary curve C of R . This is a boundary value problem and we
call it as ‘Dirichlet Problem in a Rectangle R .’ We consider a Dirichlet Problem for Laplace’s
equation (Eq. 3.146) in a Rectangle R , assuming that the temperature u ( x, y ) equals a given function
f ( x ) on the upper side and 0 on the other three sides of the rectangle. We solve this problem by
separating variables. Substituting u ( x, y ) F ( x)G ( y ) into Eq. 3.146 written as u xx u yy , dividing
1 d 2F 1 d 2G
by FG , and equating both sides to a ve constant, we get k . So, we get
F dx 2 G dy 2
d 2F d 2G
kF 0 (3.147) kG 0 (3.148)
dx 2 dy 2
For Eq. 3.147 the left and right boundary conditions imply
F (0) 0 and F ( a) 0 (3.149)
2
This gives k ( n a ) and corresponding non – zero solutions
F ( x ) Fn ( x) sin{( n a ) x} , n 1, 2, (3.150)
2
For k ( n a ) , Eq. 3.148 becomes,
2
d 2G n
2
G0 (3.151)
dy a
and solution to Eq. 3.151 is given as,
G ( y ) Gn ( y ) An e n y a Bn e n y a
(3.152)
Now the boundary condition u 0 on the lower side of R implies that Gn (0) 0 ; i.e.,
Gn (0) An Bn 0 or Bn An . This gives
Gn ( y ) An (e n y a e n y a ) 2 An sinh( n y a ) (3.153)
From Eq. 3.150 and Eq. 3.153, with 2 An An* , we thus obtain as the Eigen-functions of our problem
u n ( x, y ) Fn ( x)Gn ( y ) An* sin( n x a ) sinh( n y a ) (3.154)
These solutions satisfy the boundary condition u 0 on the left, right and the lower sides. To get a
solution also satisfying the boundary condition u ( x, b) f ( x ) on the upper side, we consider the
infinite series u ( x, y ) n 1 un ( x, y ) . From this and Eq. 3.154 with y b , we obtain
u ( x, b ) f ( x) n 1 An* sinh( n b a ) sin(n x a ) (3.155)
The expression in the parentheses of Eq. 3.155 must be the Fourier coefficients bn of f ( x ) , i.e.,
a
bn An* sinh( n b a ) (2 a ) f ( x ) sin( n x a ) dx (3.156)
0
Hence from Eq. 3.154 and 3.156, the solution of our problem is
u ( x, y ) n 1 un ( x, y ) n 1 An* sin(n x a ) sinh( n y a ) (3.157)
a
where, An* [2 {a sinh( n b a)}] f ( x) sin( n x a ) dx (3.158)
0
Statement for Linked Answer Question 3.211 and 3.212: The potential u ( x, y ) satisfies the
2u 2u
equation 0 in the square 0 x , 0 y . Three of the edges x 0 , x and
x 2 y 2
y 0 of the square are kept at zero potential and the edge y is kept at nonzero potential.
(a) u ( x, y ) n 1 An cosh nx sin ny (b) u ( x, y ) n 1 An sin nx cosh ny
(c) u ( x, y ) n 1 An sinh nx sin ny (d) u ( x, y ) n 1 An sin nx sinh ny
Solution (d): Comparing the given PDE and its conditions with Eq. 3.146 and the condition given by
Eqs. 3.149 and 3.153, we get a and b . As, the solution of Eq. 3.146, with the conditions
given by Eqs. 3.149 and 3.153, is given by Eq. 3.157. So our solution is
n x n y
un ( x, y ) n 1 An* sin sinh n 1 An* (sin nx) sinh( ny ) .
a a
Example 3.212 [XE-2007 (2 marks)]: If the edge y is kept at the potential sin x , then the
potential u ( x, y ) is given by
sin nx sinh ny sin x sinh y
(a) u ( x , y ) n 1 (b) u ( x, y )
sinh n sinh
sin x cosh y cosh nx sin ny
(c) u ( x, y ) (d) u ( x, y ) n 1
cosh cosh n
*
Solution (b): As f ( x ) sin x so the coefficient, An , in the solution given by Eq. 3.154 with n 1
and a , i.e., u ( x, y ) A* (sin nx )(sinh ny ) is given by Eq. 3.158, i.e.,
a
A* [2 {a sinh( n b a)}] f ( x) sin(n x a )dx {1 ( sinh )} 2(sin x ) sin xdx 1 {sinh(n )} .
0 0
Example 3.213 [CE-2001 (1 mark)]: The number of boundary conditions required to solve the
2 2
differential equation 2 2 0 is
x y
(a) 2 (b) 0 (c) 4 (d) 1
Solution (c): As the given equation is 2D – Heat equation, given by Eq. 3.146, so for solving it we
require four boundary conditions.
Example 3.216 [XE-2013 (1 mark)]: Which one of the following pde CANNOT be reduced to two
ordinary differentiable equations by the method of separation of variables?
u 2u 2 u 2u 2u 2u u 2 u 2 u 2u
(a) 2 0 (b) 0 (c) 0 (d) 0
t x t 2 x 2 t 2 t x x t 2 tx x 2
Solution (d): The equation given in options (a) and (b) are the heat and wave equations, respectively,
so they are solved by using method separation of variables. Now we have to check for option (c) and
u F 2u 2 F
(d). If u ( x, t ) F ( x)G (t ) is the solution then G (t ) 2 G 2 , also
x x x x
2 2 2
u u G F u G
; similarly 2
F 2 . So for option (c):
tx t x t x t t
2 u 2 u u 2 G G F F 2G F G
2
0 F 2
G ( t ) 0 F 2
G (t )
t tx x t t x x t x t
2 G t 2 1 F 2
p the variables are separated and so the equation in option (c)
(Gt ) G (t ) F x
is solved by the method of separation of variables. Now for option (d)
2u 2u 2u 2G G F 2 F
0F 2 G 2 0 the variables are not separated so the
t 2 tx x 2 t t x x
equation in option (d) is not solved by the method of separation of variables.
Example 3.218 [EC-2016 (1 mark)]: Which one of the following is a property of the solutions to the
Laplace equation: 2 f 0 ?
(a) The solutions have neither maxima nor minima anywhere except at the boundaries.
(b) The solutions are not separable in the coordinates.
(c) The solutions are not continuous.
(d) The solutions are not dependent on the boundary conditions.
Solution (a): The solutions to the Laplace equation: 2 f 0 are separable in the coordinates;
continuous; dependent on the boundary conditions; and cannot have local maxima or minima, but
extreme values must occur at the end points (i.e. at the boundaries). So option (a) is correct.
Exercise: 3.5
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
(a) u ( x, t )
2
4
m 1 4m 1
1
2
cos(2mx )e 4 m kt
2
(b) u ( x, t )
2
4
m 1 4m 1
1
2
cos(2mx)e 4 m kt
2
Copyright © 2016 by Kaushlendra Kumar e-mail: best.book4gate@gmail.com
Engineering Mathematics Chapter 3: Differential Equation [3.104]
(c) u ( x, t )
2
4
m1
4m 2 1
1
cos(2mx )e 4 m kt
2
2 4
(d) u ( x, t ) m1
1
2
4m 1 2
cos(2mx)e 4 m kt
2u 2 2u
11. The solution of PDE c , where u (0, t ) u (1, t ) 0 , u ( x, 0) A sin( x ) and
t 2 x 2
ut ( x, 0) 0 , is
(a) u ( x, t ) A sin( cx ) cos( ct ) (b) u ( x, t ) A sin( x) cos( ct )
(c) u ( x, t ) A sin( cx ) cos( t ) (d) u ( x , t ) A sin( x ) cos( t )
Answers Keys
Answer Keys: Exercise: 3.1
1 2 3 4 5 6 7 8 9 10 11 12 13 14 15
a d 1 b a b a c a c a b d c d
16 17 18 19 20 21 22 23 24 25 26 27
d 3 a d a b b c b c c b
[CS – 2017 (2 marks)]: If the ordinary generating function of a sequence {an }n 0 is (1 z ) {(1 z )3 }
, then a3 a0 is equal to _____.
3 3 ( z ) ( z)2 ( z )3
Solution: 1 (1 z ) (1 z ) 1 ( 3) ( 3)( 3 1) ( 3)( 3 1)(3 2)
1! 2! 3!
1 (1 z )3 (1 z ) 3 1 3z 6 z 2 10 z 3 15 z 4
So (1 z ) {(1 z )3 } (1 z )(1 3z 6 z 2 10 z 3 15 z 4 ) 1 4 z 9 z 2 16 z 3 25 z 4
a3 16 and a0 1 ; so a3 a0 16 1 15 .
x x2
[EC – 2017 (2 marks)]: Let f ( x ) e for real x . From among the following, choose the Taylor
series approximation of f ( x ) around x 0 , which includes all power of x less than or equal to 3.
(a) 1 x x 2 x 3 (b) 1 x (3 2) x 2 x 3 (c) 1 x (3 2) x 2 (7 6) x 3 (d) 1 x 3 x 2 7 x 3
2
Solution: The Taylor series approximation of f ( x) e x x around x0 is given by:
2 3
x0 ( x 0) ( x 0)
f ( x ) f (0) f (0) f (0) f (0)
1! 2! 3!
2 2 2 2 2
f ( x ) e x x f ( x) (1 2 x )e x x f ( x ) 2e x x (1 2 x ) 2 e x x {2 (1 2 x ) 2 }e x x
2 2
f ( x ) {4(1 2 x)}e x x {2 (1 2 x ) 2 }(1 2 x )e x x
So f (0) 1 , f (0) 1 , f (0) 3 , f (0) 7
f ( x ) 1 x (3 2) x 2 (7 6) x 3 . Thus option (c) is correct.
(a) i sin(1) (b) 2 i sin(1) (c) 3 i sin(1) (d) 4 i sin(1)
z2 1 z2 1
Solution (d): f ( z ) ez e z ; so the poles are z i , which are simple poles.
z2 1 ( z i )( z i )
( z 2 1) z ( z 2 1) z i
As Res f ( z ) lim{( z i ) f ( z )} lim ( z i) e lim e ie ; also
z i z i z i
( z i )( z i ) z i
( z i)
2
( z 1) z ( z 2 1) z i
Res f ( z ) lim{( z i) f ( z )} lim ( z i ) e lim e ie .
z i z i z i
( z i)( z i) z i
( z i)
Since z i , both poles lies inside the region | z | 3; so
2 i i
z 1 e e
I
C
z 1
2
z
e dz 2 i Res f ( z ) Res f ( z ) 2 i (ie ie ) 4 i
z i z i
i i
2i
4 i sin(1) .
[Similar question was also asked in EE – 2017 (2 marks)]
z2 1
[EE – 2017 (1 mark)]: For a complex number z , lim is
z i z 3 2 z i ( z 2 2)
(a) 2i (b) i (c) i (d) 2i
Solution (d):
z2 1 ( z i )( z i ) ( z i)( z i) ( z i)
lim 3 lim lim lim 2 2i
z i z 2 z i ( z 2 2) z i z ( z 2 2) i ( z 2 2) z i ( z i )( z 2 2) z i ( z 2)
2 2 2 2
Solution: As the curve C is y x dy dx , so I ( x i y )(dx idy ) ( x i x )(dx idx )
C C
1
2 x 1 2 x3 2 2
I (1 i) x dx (1 i) (1 i ) (1 3) (2 3)i .
x 0 3
0
z4 z3 z2
[MA – 2017 (1 mark)]: If n an ( z 2)n is the Laurent series of the function f ( z )
( z 2)3
for z \ {2} , then a2 equals _____.
1 f (z)
Solution: The coefficient an of n an ( z z0 )n is given by an
2 i C ( z z
n 1
dz ,
0)
n 0, 1, 2, . So for the given problem,
4 3 2 4 3 2
1 f (z) 1 (z z z ) 1 (z z z )
2 i C ( z 2) 2 1 2 i
C ( z 2)1 ( z 2)3 dz 2 i C
a2 dz dz , whose pole is at
( z 2) 2
z 2 , which is of 2nd order.
1 (z 4 z3 z 2 ) 1 d 21 4 3 2
2 (z z z )
Thus a2 2 i Res lim
2 1
( z 2)
2 i z 2 ( z 2) 2 (2 1)! z 2 dz ( z 2) 2
a2
1
lim
1! z 2 d
dz 1
( z 4 z 3 z 2 ) lim 4 z 3 3 z 2 2 z 32 12 4 48 .
1! z 2
2
[MA – 2017 (1 mark)]: Let C : x 2 y 2 9 be the circle in oriented positively. Then
1 2
C
{(3 y e cos x ) dx (7 x y 4 11) dy} equals _____.
Solution: Using Green’s theorem we have: C {udx vdy} R (v x) (u y ) dxdy , where R
is the region bounded by C .
x2
Now let C {(3 y e cos )dx ( 7 x y 4 11)dy}
2
So u ( x, y ) (3 y e cos x ) u y 3 ; v ( x, y ) ( 7 x y 4 11) vx 7 ;
[PE – 2017 (1 mark)]: If 5 x 2iy ix 7 y 2 3i , where i 1 , the values of two real numbers
( x, y ) are, respectively:
(a) ( 1,1) (b) (1, 1) (c) (1,1) (d) ( 1, 1)
Solution (a): 5 x 2iy ix 7 y 2 3i (5 x 7 y ) i(2 y x) 2 3i 5 x 7 y 2 …(i) and
2 y x 3 …(ii). From (i) and (ii), we have x 1 , y 1 .
[PE – 2017 (1 mark)]: Pick the INCORRECT inequality, where z1 , z2 and z3 are complex numbers.
(a) | z1 z2 | | z1 | | z2 | (b) | z1 z2 | | z1 | | z2 |
(c) | z1 z2 | | z1 | | z2 | (d) | z1 z2 z3 | | z1 | | z2 | | z3 |
i i i
Solution (c): Let z1 r1e z 2 r2 e , z3 r3e
z1 z 2 r1e i r2 ei r1 (cos i sin ) r2 (cos i sin ) ( r1 cos r2 cos ) i ( r1 sin r2 sin )
So z1 z 2 (r1 cos r2 cos ) i ( r1 sin r2 sin ) ( r1 cos r2 cos ) 2 (r1 sin r2 sin ) 2
dz
[PH – 2017 (1 mark)]: The contour integral 1 z 2 evaluated along a contour going from to
along the real axis and closed in the lower half-plane by a half circle is equal to _____. (up to two
decimal places).
1
Solution: f ( z ) , whose poles are z i and z i , of which both are simple poles. But as
1 z2
shown in figure z i is not lying; but z i is lying in the contour C . So
dz
1 z 2 2 i Resz i
f ( z ) 2 i lim( z i) f ( z )
z i
dz 1 1
1 z 2 2 i lim(
z i
z i)
( z i)( z i)
2 i lim
( z i)
z i
dz
1 z 2 2 i 1
2i
3.14 .
u v
Also f ( y ) 2 y f ( y ) y 2 c …(ii)
y x
From (i) and (ii), we get: u ( x, y ) x 2 3 x y 2 c . As u (0, 0) 0 c 0 .
Thus u ( x, y ) x 2 3x y 2 . So u (1,1) 12 3 1 12 3 .
[PI – 2017 (1 mark)]: The product of a complex number z x iy and its conjugate z is
(a) x 2 (b) y 2 (c) x 2 y 2 (d) x 2 y 2
Solution (d): zz ( x iy )( x iy ) x 2 i 2 y 2 x 2 y 2
[XE – 2017 (1 mark)]: Let C be a simple smooth closed curve enclosing the region R in the xy
plane. Let C be oriented counter-clockwise. If the value of the integral
x2
C {( y e ) dx (3x cos y ) dy} is 16, then the area of R is _____.
Solution: Using Green’s theorem we have: C {udx vdy} R (v x) (u y ) dxdy , where R
is the region bounded by C .
x2 x2
Now let C {( y e ) dx (3 x cos y ) dy}
{( y e )dx (3x cos y)dy}
C
2
So u ( x, y ) y e x u y 1 ; v ( x, y ) 3x cos y vx 3 ;
Hence C {udx vdy} R (3 1)dxdy 2R dxdy 16 R dxdy 8 .
1 1 i 1 ei t 2 1 1
[XE – 2017 (2 marks)]: Let I dt i t 2
dt i dt , where t is a real variable
0 1 t 0
2 1 e 0 1 it
and i 1 . The value of I is _____.
i 1 ei t 2
1 1 1 1 1 i t 2 1 1
Solution: I dt
i t 2
dt i dt {ln(1 t )}0 {ln(1 e )}0 {ln(1 it )}0
0 1 t 0
2 1 e 0 1 it
i 2
I ln 2 ln(1 e ) ln 2 ln(1 i) ln{1 cos( 2) i sin( 2)} ln(1 i) 0 .
Real and Imaginary part of a Complex Number: If x and y are two real numbers, then a
number of the form z x iy is called a complex number. Here ‘ x ’ is called the real part of z
[denoted by Re( z ) ] and ‘ y ’ is known as the imaginary part of z [denoted by Im( z ) ].
A complex number z is purely real if its imaginary part is zero i.e., Im( z ) 0 and purely
imaginary if its real part is zero i.e., Re( z ) 0 .
A real number can be written as a i.0 , therefore every real number can be considered as a
complex number whose imaginary part is zero.
Complex number as an ordered pair: A complex number ‘ a ib ’ may also be defined as an
ordered pair of real numbers and denoted by ( a, b) . For e.g. i can be denoted by (0,1) .
The algebraic operations performed on two complex numbers z1 x1 iy1 and z2 x2 iy2 are:
Addition/Subtraction: ( x1 iy1 ) ( x2 iy2 ) ( x1 x2 ) i( y1 y2 )
Multiplication: ( x1 iy1 )( x2 iy2 ) ( x1 x2 y1 y2 ) i ( x1 y2 y1 x2 )
x1 iy1 x1 iy1 x2 iy2 ( x1 x2 y1 y2 ) i ( y1 x2 y2 x1 )
Division: 2 2
2 2
, where at least one of
x2 iy2 x2 iy2 x2 iy2 x2 y2 x2 y2
x2 and y2 is non-zero.
Equality of two Complex Numbers: Two complex numbers z1 x1 iy1 and z2 x2 iy2 are
said to be equal if and only if their real parts and imaginary parts are separately equal. i.e.
z1 z2 x1 iy1 x2 iy2 x1 x2 and y1 y2 . Thus, one complex equation is equivalent to
two real equations. A complex number z x iy 0 iff x 0, y 0 .
The complex number do not possess the property of order i.e., ( a ib) (or) (c id ) is not
defined. For example, the statement 9 6i 3 2i makes no sense.
Reciprocal of a complex number: For an existing non-zero complex number z x iy , the
1 1 z 1 1 x iy Re( z ) i[ Im( z )] z
reciprocal is given by z 2
, i.e. z 2 2
2
2
2
.
z |z| x iy x y |z| |z| |z|
Example 4.3 [ME-2011 (1 mark)]: The product of two complex numbers 1 i and 2 5i is
(a) 7 3i (b) 3 4i (c) 3 4i (d) 7 3i
2
Solution (a): (1 i )(2 5i ) 2 5i 2i 5i 2 3i 5 7 3i .
[Similar question was also asked in PI-2009 (1 mark)]
Example 4.5 [IN-2008 (2 marks)]: A complex variable Z x j 0.1 has its real part x varying in
the range to . Which one of the following is the locus (shown in thick lines) of 1 Z in the
complex plane?
1 Z Z x j 0.1 x 0.1
Solution (b): Y 2
2 2
Re(Y ) 2 2
and Im(Y ) . As
Z ZZ Z x 0.1 x 0.1 x 0.12
2
Im(Y ) Re(Y ) 0.1 x Re(Y ) 10 x Im(Y ) which means the slope of Y keep changing as
x and so option (a) and (d) are not correct. As
1
{Re(Y )}2 {Im(Y )}2 2 2
10 Im(Y ) {Re(Y )}2 {Im(Y )}2 10 Im(Y ) 0
x 0.1
{Re(Y )} {Im(Y ) 5}2 52 . So locus of Y is a circle whose centre is at (0, 5) and radius 5,
2
these two conditions are satisfied by the figure given in option (b). So option (b) is correct.
Example 4.6 [EC-2005 (2 marks)]: In what range should Re( s) remain so that the Laplace
transform of the function e( a 2) t 5 exists?
(a) Re( s ) a 2 (b) Re( s ) a 7 (c) Re( s) 2 (d) Re( s ) a 5
Solution (a): We have f (t ) e ( a 2) t 5 e5 e ( a 2) t { f (t )} {e5 e ( a 2) t } F ( s) e5 {e ( a 2) t }
F ( s ) e5 [1 {s ( a 2)}] . Now as Laplace transform of any function is always positive so
s ( a 2) 0 s a 2 ; as s is the Re( s) so we have Re( s ) a 2 .
Example 4.8 [CH-2014 (1 mark)]: If f * ( x ) is the complex conjugate of f ( x) cos x i sin x , then
b
for real a and b , a f * ( x ) f ( x ) dx is always
(a) Positive (b) Negative (c) real (d) imaginary
Solution (c): The complex conjugate of f ( x ) cos x i sin x e is f * ( x ) cos x i sin x e ix , so
ix
b b
f * ( x ) f ( x ) e ix eix 1 thus I f * ( x ) f ( x )dx 1dx I b a which can be positive or
a a
negative but it is definitely real as a and b are real.
Example 4.9 [EE-2014 (1 mark)]: Let S be the set of points in the complex plane corresponding to
the unit circle. (That is S z : z 1 ). Consider the function f ( z ) zz * where z * denotes the
complex conjugate of z . The f ( z ) maps S to which one of the following in the complex plane
(a) unit circle (b) horizontal axis line segment from origin to (1, 0)
(c) the point (1, 0) (d) the entire horizontal axis
2
Solution (c): As z* z , the conjugate of z ; thus f ( z ) z z z 1 1 i (0) f ( z ) is a point
(1, 0) in the complex plane.
Example 4.11 [CE-2005 (1 mark)]: Which of the following is not true for complex number Z1 , Z 2 ?
2
(a) Z1 Z 2 ( Z1 Z2 ) Z 2 (b) Z1 Z 2 Z1 Z 2
2 2 2 2
(c) Z1 Z 2 Z1 Z 2 (d) Z1 Z 2 Z1 Z 2 2 Z1 2 Z 2
2 2
Solution (b): For option (a): As Z Z Z so RHS ( Z1 Z 2 ) Z 2 ( Z1 Z 2 ) ( Z 2 Z 2 ) Z1 Z 2 LHS .
For option (b) and (c), refer to triangle inequality given by Eq. 4.2, so we can say that option (b) is not
2 2
correct. For option (d): Let S Z1 Z 2 Z1 Z 2 so S ( Z1 Z 2 )( Z1 Z 2 ) ( Z1 Z 2 )( Z1 Z 2 )
S ( Z1 Z 2 )( Z1 Z 2 ) ( Z1 Z 2 )( Z1 Z 2 )
S Z1 Z1 Z1 Z 2 Z 2 Z1 Z 2 Z 2 Z1 Z1 Z1 Z 2 Z 2 Z1 Z 2 Z 2
2 2 2 2 2 2
S Z1 Z 1 Z1 Z 2 2 Z1 2 Z 2 .
Example 4.12 [CH-2016 (1 mark)]: What are the modulus (r ) and argument ( ) of the complex
number 3 4i ?
(a) r 7 , tan 1 (4 3) (b) r 7 , tan 1 (3 4)
(c) r 5 , tan 1 (3 4) (d) r 5 , tan 1 (4 3)
Solution (d): r 32 4 2 5 and tan 4 3 tan 1 (4 3) .
Multiplication and Division in Polar Form: Let z1 x1 iy1 r1 (cos 1 i sin 1 ) and
z2 x2 iy2 r2 (cos 2 i sin 2 ) , then
Multiplication: z1 z2 r1r2 (cos 1 cos 2 sin 1 sin 2 ) i (sin 1 cos 2 cos 1 sin 2
z1 z2 r1r2 cos(1 2 ) i sin(1 2 ) (4.4)
z1 z2 r1r2 r1 r2 z1 z2 (4.5)
also, arg( z1 z2 ) arg z1 arg z2 (upto multiples of 2 ) (4.6)
Division: z1 ( z1 z2 ) z2 z1 ( z1 z2 ) z2 . Hence by product rule of two complex numbers,
z1 ( z1 z2 ) z2 z1 z2 z1 z2 (4.7)
Also, by product rule, we have arg( z1 ) arg ( z1 z2 ) z2 arg( z1 z2 ) arg( z2 )
arg( z1 z2 ) arg( z1 ) arg( z2 ) (upto multiples of 2 ) (4.8)
z1 z2 (r1 r2 ) cos(1 2 ) i sin(1 2 ) (4.9)
Example 4.13 [ME-2010 (1 mark)]: The modulus of the complex number (3 4i) (1 2i) is
(a) 5 (b) 5 (c) 1 5 (d) 1 5
z1 z1 3 4i 3 4i 32 4 2 5
Solution (b): As 5
z2 z2 1 2i 1 2i 2 2
1 ( 2) 5
[Similar question was also asked in CH-2009 (1 mark)]
Example 4.15 [CS-1995 (2 marks)]: If the cube roots of unity are 1, and 2 , then the roots of the
following equation are ( x 1) 3 8 0
2 2 2 2
(a) 1,1 2 ,1 2 (b) 1,1 2 ,1 2 (c) 1,1 2 ,1 2 (d) 1,1 2, 1 2
2
Solution (c): As x 1 ( 8)1 3 ( 2)3 3 (1)1 3 2(1)1 3 2, 2 , 2 2 x 1,1 2 ,1 2 .
Example 4.16 [IN-2009 (2 marks)]: One of the root of the equation x 3 i , where i is positive
square root of 1 , is
(a) i (b) ( 3 2) i(1 2) (c) ( 3 2) i (1 2) (d) ( 3 2) i(1 2)
Solution (b): x 3 i i (i 2 ) i 3 x 3 (i )3 (1)3 x i (13 )1 3 x i (1)1 3 . As cube root of
unity are 1, , 2 , where {( 1 2) i( 3 2)} and 2 {( 1 2) i ( 3 2)} . So
1 3 1 3 i 3 i 3
x i , i i , i i i , , . Thus from the given options we
2 2 2 2 2 2 2 2
can say that option (b) is correct.
Example 4.17 [PI-2010 (1 mark)]: If a complex number satisfies 3 1 , then the value of
1 (1 ) is
(a) 0 (b) 1 (c) 2 (d) 4
Solution (a): If x 3 1 then its roots are 1, {( 1 2) i( 3 2)} and 2 {( 1 2) i ( 3 2)}
which satisfies 1 2 0 and 3 1 . So 1 (1 ) ( 2 1) 0 .
Example 4.18 [EE-2011 (2 marks)]: The two vectors 1 1 1 and 1 a a 2 , where
a (1 2) i( 3 2) are
(a) Orthonormal (b) Orthogonal (c) Parallel (d) Collinear
Solution (b): As cube roots of unity are 1, {( 1 2) i( 3 2)} and 2 {( 1 2) i ( 3 2)}
which satisfies 1 2 0 and 3 1 . Now as the dot product of the given vectors is
2 T 2
1 1 1 1 a a 1 a a 0 , as we have a , so the given vectors are orthogonal, i.e.,
perpendicular to each other and thus they are neither parallel nor collinear. Also both the vectors are
not of unit length so they are not orthonormal.
the values of f are tending to l for all z tending to z0 . In other words, if ve real we can find
a ve real s.t. z z0 in the disk z z0 , then f ( z ) l . If a limit exists, it is unique.
f ( z 0 z ) f ( z 0 )
the point z0 is denoted by f ( z0 ) and is defined as f ( z0 ) lim provided this
z z 0
limit exists. The first thing that we notice is that z , being a complex number, can approach zero in
more than one way. If we write z x iy , then we can approach zero along the real axis y 0
or along the imaginary axis x 0 , or indeed along any direction. For the derivative to exist, the
answer should not depend on how z tend to zero. Now let us see what this means: We have
f ( z0 ) u ( x0 , y0 ) iv ( x0 , y0 ) and f ( z0 z ) u ( x0 x, y0 y ) iv ( x0 x, y0 y ) , so
u ( x0 , y0 ) iv ( x0 , y0 )
f ( z0 ) lim where u ( x0 , y0 ) u ( x0 x, y0 y ) u ( x0 , y0 ) ,
x 0, y 0 x iy
v( x0 , y0 ) v ( x0 x, y0 y ) v( x0 , y0 ) . Let us first take the limit z 0 by first taking y 0
and then x 0 , in other words we let z 0 along the real axis, then we have
u ( x0 , y0 ) iv ( x0 , y0 ) u ( x0 , y0 ) iv( x0 , y0 ) u v
f ( z0 ) lim lim lim i
x 0 y 0 x iy x 0 x x ( x0 , y0 ) x ( x0 , y0 )
Now let us take the limit z 0 by first taking x 0 and then y 0 , in other words we let
z 0 along the imaginary axis, then we have
u ( x0 , y0 ) iv ( x0 , y0 ) u ( x0 , y0 ) iv ( x0 , y0 )
f ( z0 ) lim lim lim
y 0 x 0 x iy x 0 iy
iu ( x0 , y0 ) v ( x0 , y0 ) u v
f ( z0 ) lim i .
x 0 y y ( x0 , y0 ) y ( x0 , y0 )
The differentiation rules are the same as in real calculus as their proofs are literally the same.
Thus for any analytic functions f and g and constant c we have,
(cf ) cf , ( f g ) f g , ( fg ) f g f g , ( f g ) ( f g f g ) g 2 .
hence it is well defined for all z0 . On the other hand, the function f ( z ) z , where z is the conjugate
f ( z0 z ) f ( z0 ) z0 z z 0 z
of z , is not differentiable anywhere as f ( z0 ) lim lim lim
z z 0 z zz 0 z 0
; now if we let z tend to zero along real values, we would find that f ( z0 ) 1 , whereas if we would
let z tend to zero along imaginary values we would find that f ( z0 ) 1 . We could have reached
the same conclusion via the Cauchy-Riemann equations with u ( x, y ) x and v ( x, y ) y which
violates the Cauchy-Riemann equation.
The analyticity, unlike differentiability, is not a property of a function at a point, but on an open
set of points. The reason for this is to eliminate from the class of interesting functions, functions
which may be differentiable at a point but nowhere else. Whereas this is a rarity in calculus, it is a
very common occurrence for complex-valued functions of a complex variables. For example,
2
consider the function f ( z ) z , this function has u ( x, y ) x 2 y 2 and v ( x, y ) 0 . Therefore
the Cauchy-Riemann equations are only satisfied at the origin in the complex plane because:
u x 2 x v y 0 x 0 and vx 0 u y 2 y y 0 .
v y v x 3( x y ) 2 2kxy ky 2 …(iv) [by applying (i)]; also partial differentiation of (iii) w.r.t. y ,
we get u y v y 3( x y ) 2 kx ( x y ) kxy ( 1) 3( x y ) 2 2kxy kx 2
vx v y 3( x y ) 2 2kxy kx 2 …(v). Now (iv) + (v) 2v x ky 2 kx 2 …(vi); also (iv) –
(v) 2v y 6( x y ) 2 4kxy ky 2 kx 2 …(vii). Now integrating (vi) w.r.t. x we get
2v ky 2 x k ( x 3 3) f ( y ) …(viii); also integrating (vii) w.r.t. y we get
3 2 3 2
2v 2( x y ) 2kxy k ( y 3) kx y h( x ) …(ix). Thus (viii) + (ix)
2 3 3 2 3 2
ky x k ( x 3) f ( y ) 2( x y ) 2kxy k ( y 3) kx y h( x ) 0
xy 2 ( k 6) x 2 y (6 k ) (k 3) 2 x 3 2 ( k 3) y 3 f ( y ) h ( x ) 0 . As f ( y ) and h( x) both
are pure function of y and x , respectively so they do not contain the mixed terms of x and y . Thus
the coefficient of xy 2 and x 2 y must be zero and so we have k 6 0 and 6 k 0 k 6 .
Example 4.27 [EC-2014 (2 marks)]: The real part of an analytic function f ( z ) where z x jy is
given by e y cos x . The imaginary part of f ( z ) is
(a) e y cos x (b) e y sin x (c) e y sin x (d) e y sin x
Solution (b): If f ( x iy ) u ( x, y ) iv( x, y ) is analytic then it satisfies the Cauchy-Riemann
equations, i.e., (i) u x v y , (ii) u y vx . So (i) (v y ) ( x )(e y cos x) (v y ) e y sin x
v e y sin x y , on integrating both sides we get, v e
y
sin x y v e y sin x f ( x )
…(iii). Now differentiating (iii) both sides w.r.t. x and applying the (ii) condition we get
v
{e y cos x f ( x )} (e y cos x ) e y cos x f ( x ) e y cos x f ( x) 0
x y
f ( x ) k , where k is any constant. Thus from (iii) we have v e y sin x k .
[Similar questions were also asked in ME-2009, CE-2011, ME-2014, ME-2014 (2 marks), PI-
2010 (1 mark)]
Example 4.35 [PE-2016 (1 mark)]: For a complex number Z (1 2) ( 3 2)i , the value of Z is 6
(a) (1 2) ( 3 2)i (b) 1 (c) (1 2) ( 3 2)i (d) 1
i
Solution (d): Converting the given complex number, z x iy , in polar form we get z re , where
r x2 y 2 and tan 1 ( y x ) . So we have x 1 2 and y 3 2, thus
Complex Trigonometric Functions: Equation 4.16 implies that ei (cos i sin ) and
e i e i ( ) (cos i sin ) . By adding and subtracting these two equations we get
cos (e i e i ) 2 and sin (ei e i ) (2i) (4.17)
Hence Equations 4.17 suggests the following definitions for complex values z x iy as,
cos z (eiz e iz ) 2 and sin z (eiz e iz ) (2i) (4.18)
sin z cos z 1 1
Also, tan z , cot z , sec z and cosec z (4.19)
cos z sin z cos z sin z
Since e z is entire, cos z and sin z are entire functions. tan z and sec z are not entire; they are
analytic except at the points where cos z is zero; cot z and cosec z are analytic except where sin z is
zero. All the formulas for derivative of real trigonometric function are also applied to complex
trigonometric functions, i.e. (cos z ) sin z , (sin z ) cos z , (tan z ) sec 2 z .
Complex Hyperbolic Functions: The complex hyperbolic cosine and sine functions are given as,
cosh z (e z e z ) 2 and sinh z (e z e z ) 2 (4.20)
These functions are entire with derivatives (cosh z ) sinh z and (sinh z ) cosh z as in calculus.
The other hyperbolic functions are defines as,
sinh z cosh z 1 1
tanh z , coth z , sec hz and cosec hz (4.21)
cosh z sinh z cosh z sinh z
Example 4.37 [CH-2012 (1 mark)]: If a and b are arbitrary constants, then the solution to the
ordinary differential equation d 2 y dx 2 4 y 0 is
(a) y ax b (b) y ae x (c) y a sin 2 x b cos 2 x (d) y a cosh 2 x b sinh 2 x
Solution (a): As we have homogeneous linear equation with constant coefficient so putting
y e mx y me mx y m 2 e mx in the given differential equation we get e mt ( m 2 4) 0 so we
have the auxiliary equation m 2 4 0 as e mt 0 . As the discriminant of the auxiliary equation
D 0 2 4 1 ( 4) 16 0 so we have distinct real root case; also the roots of our auxiliary
equation are m1 , m2 2, 2 ; so the general solution of given differential equation is
y c1e 2 x c2 e 2 x . Now from Eq. 4.28 we have e z cosh z sinh z so e 2 x cosh(2 x) sinh(2 x) and
e 2 x cosh( 2 x) sinh( 2 x) cosh(2 x) sinh(2 x) . Thus
2 x 2x
y c1e c2 e c1{cosh(2 x ) sinh(2 x )} c2 {cosh(2 x ) sinh(2 x)}
y (c1 c2 ) cosh(2 x ) (c1 c2 ) sinh(2 x) a cosh(2 x) b sinh(2 x) , a (c1 c2 ) and b (c1 c2 ) .
Example 4.39 [IN-2009 (1 mark)]: If z x iy , where x and y are real. The value of eiz is
(a) 1 x2 y 2 (c) e y (d) e y
(b) e
Solution (d): As eiz ei ( x iy ) e y ix e y eix e y (cos x i sin x ) , so
iz y y y 2 2 y
e e (cos x i sin x) e cos x i sin x e cos x sin x e .
y cos i sin y cos i sin , cos i sin
4 4 4 4 4 4
y cos( 4) i sin( 4), cos(3 4) i sin(3 4) , as cos( 4) cos(3 4) and
sin( 4) sin(3 4) .
Example 4.41 [EC-2006 (2 marks)]: For the function of a complex variable W ln Z (where,
W u iv and Z x iy ), the u constant lines get mapped in z plane as
(a) set of radial straight lines (b) set of concentric circles
(c) set of confocal hyperbolas (d) set of confocal ellipses
Solution (b): As u iv ln( x iy ) so from Eq. 4.32 we have ln( x iy ) ln r i , where
r x2 y 2 and tan 1 ( y x ) ; so ln( x iy ) ln x 2 y 2 i tan 1 ( y x) . Thus
u iv ln x 2 y 2 i tan 1 ( y x ) u ln x 2 y 2 x 2 y 2 e 2u . It is given that u is a constant
so e 2u is also constant, say, k . So x 2 y 2 k which is a set of concentric circles.
Example 4.42 [IN-2013 (1 mark)]: The complex function tanh( s) is analytic over a region of the
imaginary axis of the complex s plane if the following is TRUE everywhere in the region for all
integers n
Exercise: 4.1
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
2 200
1. If i 1, then the value of n1 i n _____.
2. If {(1 i ) (1 i)}x 1 then
(a) x 4 n , where n is any positive integer (b) x 2 n , where n is any positive integer
(c) x 4 n 1 , where n is any positive integer (d) x 2 n 1 , where n is any positive integer
3. If x 3 i , then x 3 3x 2 8 x 30 _____
4. 1 (1 2i) 3 (1 i ) (3 4i ) (2 4i )
(a) (1 2) (9 2)i (b) (1 2) (9 2)i (c) (1 4) (9 4)i (d) (1 4) (9 4)i
4 2
5. The real values of x and y for which ( x 2 xi ) (3 x yi) (3 5i) (1 2 yi ) is satisfied, are
(a) x 2, y 3 (b) x 2, y 1 3 (c) Both (a) and (b) (d) None of these
6. The complex numbers sin x i cos(2 x ) and cos x i sin(2 x ) are conjugate to each other for
(a) x n (b) x n (1 2) (c) x 0 (d) No value of x
1
7. The real part of (1 cos 2i sin ) is
(a) 1 (3 5 cos ) (b) 1 (5 3 cos ) (c) 1 (3 5 cos ) (d) 1 (5 3cos )
8. If and are different complex numbers with 1 , then ( ) (1 ) _____.
9. Amplitude of {(1 i ) (1 i )} is
(a) 2 (b) 2 (c) 4 (d) 6
10. Let z , w be complex numbers such that z iw 0 and arg( zw) , then arg( z )
(a) 5 4 (b) 2 (c) 3 4 (d) 4
11. If z and are to non-zero complex numbers such that z 1 and arg( z ) arg( ) 2 then
z is equal to
(a) 1 (b) 1 (c) i (d) i
12. The square root of 3 4i are
(a) (2 i) (b) (2 i) (c) ( 3 2i ) (d) ( 3 2i )
13. 2i equals
(a) 1 i (b) 1 i (c) 2i (d) None of these
2
14. (1 7i ) (2 i )
3 3 3 3 (d) None of
(a) 2 cos i sin (b) 2 cos i sin (c) cos i sin
4 4 4 4 4 4 these
i
15. Real part of e e is
(a) e cos [cos(sin )] (b) e cos [cos(cos )] (c) esin [sin(cos )] (d) esin [sin(sin )]
16. i log{( x i ) ( x i)} is equal to
(a) 2 tan 1 x (b) 2 tan 1 x (c) 2 tan 1 x (d) 2 tan 1 x
Existence of the complex line integral: Our assumptions that f ( z ) is continuous and C is
piecewise smooth imply the existence of the line integral Eq. 4.27. Let f ( z ) u ( x, y ) iv ( x , y ) ,
z (t ) x (t ) iy (t ) dz dx (t ) idy (t ) , a t b . Then Equation 4.27 can be written as,
We can write z0 and z1 instead of C , as we get same value for all those C from z0 to z1 .
Indefinite integration of any continuous complex function: This method is not restricted to
analytic functions but applies to any continuous complex function. Let C be a piecewise smooth
path, represented by z z (t ) , where a t b . Let f ( z ) be a continuous function on C . Then,
b
C f ( z )dz a f z (t ) z(t )dt , where ( z dz dt ) (4.31)
Proof: The left hand side of Eq. 4.31 is given by Eq. 4.29 in terms of real line integrals, and we show
that the right hand side of Eq. 4.31 also equals to Eq. 4.29. We have z x iy z x iy . We
and dy ydt
simply write u for u x(t ), y (t ) and v for v x (t ), y (t ) . We also have dx xdt .
Consequently, in Eq. 4.31 we have
b b
a f z (t ) z(t )dt a (u iv)( x iy )dt C udx vdy i(udy vdx) C (udx vdy ) i C (udy vdx)
If we integrate a given function f ( z ) from a point z0 to a point z1 along different paths, the
integrals will in general have different values. In other words, a complex line integral depends not
only on the endpoints of the path but in general also on the path itself.
Bound for Integrals: The absolute value of complex line integrals is given by
C f ( z )dz ML (4.32)
where, L is the length of C and M a constant such that f ( z ) M everywhere on C .
A simple closed path is a closed path that does not intersect or touch itself. For example, a circle
is simple, but a curve shaped like an 8 is not simple, as shown in Fig. 4.7(a) and (b).
A simply connected domain D in
the complex plane is a domain such
that every simple closed path in D
encloses only points of D , for e.g.,
the interior of a circle (‘open disk’),
ellipse, or any simple closed curve,
as shown in Fig. 4.7(c). A domain
that is not simply connected is called
multiply connected; for e.g., an Figure 4.7: (a), (b) are Closed Paths; (c), (d) and (e) are simply and
annulus, a disk without the centre, as multiply connected domains
shown in Fig. 4.7(d) and (e).
Proof: From Eq. 4.29, we have C f ( z)dz C (udx vdy) i C (udy vdx) . Since f ( z ) is analytic
in D , its derivative f ( z ) exists in D . Since f ( z ) is assumed to be
continuous and u and v have continuous partial derivatives in D . Hence
Green’s theorem (with u and v instead of F1 and F2 ) is applicable and
gives C (udx vdy ) R (v x) (u y) dxdy , where R is the region
bounded by C . Figure 4.8: Cauchy –
Integral Theorem
The second Cauchy – Riemann shows that the integrand on the right is identically zero. Hence the
integral on the left is zero. Also by use of first Cauchy – Riemann equation that the last integral in the
above formula is zero.
Example 4.43 [IN-2009 (1 mark)]: The value of (1 z ) sin zdz , where (a) 0 (b) 2 i
the contour of integration is a simple closed curve around the origin, is (c) (d) 1 2 i
sin z 1 z3 z5 z 7 z2 z4 z6
Solution (a): As the function f ( z ) z 1 which
z z 3! 5! 7! 3! 5! 7!
is analytic in a simple closed curve around the origin so from Eq. 4.43 we have
(1 z ) sin zdz 0 .
Cauchy’s Integral Theorem 2: If f ( z ) is analytic in a simply connected domain D , then the
integral of f ( z ) is independent of path in D .
Proof: Let z1 and z2 be any points in D . Consider two paths C1 and C 2 in D from z1 and z2
without further common points, as shown in Fig. 4.9a. Denote by C 2* the path C 2 with the orientation
reserved, as shown in Fig. 4.9b. Integrate from z1 over C1 to z2 over C 2* back to z1 . This is a simple
closed path and Cauchy theorem applies under our assumption of the present theorem and gives zero,
i.e. fdz * fdz 0 fdz * fdz . But the minus sign on the right disappear if we integrate
C1 C2 C1 C2
in the reverse direction, from z1 to z2 which shows that the integral of f ( z ) over C1 and C 2 are
equal, i.e. C1
f ( z )dz
C2
f ( z ) dz . This proves the theorem for paths that have only the endpoints in
common. For paths that have
finitely many further common
points, apply the present
argument to each ‘loop’
(portion of C1 and C 2
between consecutive common
points; four loops as shown in Figure 4.9: (a) for path function only end points are in common (b) Cauchy integral
Fig. 4.9c). for reversed orientation (c) Paths with more common points
3i
Example 4.44 [ME-2014 (2 marks)]: If z is a complex variable, the value of 5 (1 z ) dz is
(a) 0.511 1.57i (b) 0.511 1.57i (c) 0.511 1.57i (d) 0.511 1.57i
Solution (b): As f ( z ) 1 z is analytic between the line joining the points z 5 and z 3i .
3i
So I (1 z )dz (ln z )35i ln 3i ln 5 ln 3 ln i ln 5 ln(3 5) ln e i 2
0.5108 i( 2) (as
5
Example 4.45 [MA-2014 (2 marks)]: Let {z : Im( z ) 0} and let C be a smooth curve lying
in with initial point 1 2i and final point 1 2i . The value of C {(1 2 z) (1 z )}dz is
1 1 1 1
(a) 4 ln 2 i (b) 4 ln 2 i (c) 4 ln 2 i (d) 4 ln 2 i
2 4 2 4 2 4 2 2
Solution (a): As it is given that f ( z ) (1 2 z ) (1 z ) is a smooth curve lying in with initial point
1 2i and final point 1 2i . So
1 2z 1 2 i 1 2 z 1 2 i 1 2 z 1 1 1 2i 2 2 z 1 1 2i 1
I dz dz dz dz 2 dz
C 1 z 1 2i 1 z 1 2i 1 z 1 2 i 1 z 1 2 i
1 z
1 2 i
I 2 z ln(1 z )1 2i 2{(1 2i ) (1 2i )} {ln(1 1 2i) ln(1 1 2i)}
1 i (1 i )i 1 ( 1)
I 4 {ln(2 2i) ln(2i)} 4 ln 4 ln 2
4 ln(1 i ) 4 ln{ 2e i tan }
i i
i 4
I 4 ln 2 ln{e } 4 (1 2) ln 2 i( 4)
Cauchy integral theorem for multiply connected domains: Cauchy’s theorem applies to
multiply connected domains. We first explain this for a doubly connected domain D with outer
boundary curve C1 and inner C 2 as shown in Fig. 4.10(a). If a function f ( z ) is analytic in any
domain D * that contains D and its boundary curves then we have,
f ( z )dz f ( z )dz
C1 C2
(4.34)
both the integrals being taken counter-clockwise (or both clockwise, and regardless of whether or not
the full interior of C 2 belongs to D * )
Proof: By two cuts A and B , as shown in Fug. 4.10(b), we cut D into two simply connected
domains D1 and D2 in which and on whose boundaries f ( z ) is analytic. By Cauchy’s integral
theorem the integral over the entire boundary of D1 (taken in the sense of arrows, as shown in Fig.
4.10(b), is zero, and so is the integral over the boundary of D2 , and thus their sum. In this sum the
integrals over the cuts C1 and C 2 cancel because we integrate over them in both directions and we
are left with the integrals over C1 (counter-clockwise) and C 2 (clockwise); hence by reversing the
integration over C 2 we have, C
1
fdz
C2
fdz 0 , which is same as Eq. 4.34. For domains of higher
connectivity the idea remains the same. Thus, for a triply connected domain we use three cuts A, B, C
as shown in Fig. 4.10(c).
Adding integrals as
before, the integrals over
the cuts cancel and the
sum of the integrals over
C1 (counter-clockwise)
and C 2 , C3 (clockwise) Figure 4.10: (a), (b) are Doubly connected domain; (c) is Triply connected domain
is zero. Hence the integral over C1 equals the sum of the integrals over C 2 and C3 , all three now
taken counter-clockwise. Similarly for quadruply connected domains and so on.
Cauchy Integral Formula: Let f ( z ) be analytic in a simply connected domain D . Then for any
point z0 in D and any simple closed path C in D that encloses z0 as shown in Fig. 4.11(a), then
f ( z)
C z z dz 2 if ( z0 ) (4.35)
0
the integration being taken counter-clockwise. For
representing f ( z0 ) by a contour integral, divide Eq. 4.35 by
2 i , we get
1 f (z) Figure 4.11: (a) Cauchy integral formula (b)
f ( z0 )
2 i C z z0
dz (4.36) Multiply connected domain
For multiply connected domains, if f ( z ) be analytic on C1 and C 2 and in the ring – shaped domain
bounded by C1 and C 2 , as shown in Fig. 4.11(b), and z0 is any point in that domain, then (where the
outer integral (over C1 ) is taken a clockwise and the inner integral (over C 2 ) is taken clockwise)
1 f ( z) 1 f ( z)
f ( z0 )
2 i C 1 z z0
dz
2 i C2 z z0
dz (4.37)
2! f ( z)
2 i C ( z z )3
f ( z0 ) dz (4.39)
0
Figure 4.12: Derivative of an analytic
n! f ( z) function
2 i
( n)
In general, f ( z0 ) C ( z z0 )n 1 dz (4.40)
here C is any simple closed path in D that encloses z0 and whose full interior belongs to D ;
and we integrate counter-clockwise around C , as shown in Fig. 4.12.
Cauchy’s inequality: Let for Eq. 4.40, we choose C as a circle of radius r and centre z0 and
apply the ML inequality (Eq. 4.42) with f ( z ) M on C we obtain from Eq. 4.40 as,
n! f (z) n! 1 n!M
C ( z z0 )n1 dz 2 M r n1 2 r
(n) (n)
f ( z0 ) n
(4.41) f ( z0 )
2 r
Liouville’s theorem: It states that, if an entire function is bounded in absolute value in the whole
complex plane, then this function must be a constant.
If f ( z ) is continuous in a simply connected domain D and if C f ( z )dz 0 for every closed
path in D , then f ( z ) is analytic in D .
Example 4.46 [XE-2011 (1 Solution (c): As f ( z ) z 3e z and all its derivatives are analytic in
mark)]: The integral
given curve C . So using Eq. 4.39, we have
C
3 z 3
z e ( z 1) dz along the
2! z 3e z z 3e z
curve C : z 2 , oriented
f ( z 0 )
2 i C ( z 1)3
dz C ( z 1)3 dz i f (1) . As
Tests for Convergence and Divergence of Series: The tests are practically the same as in
calculus. We apply them before we use a series, to make sure that the series converges.
Divergence Test: If a series z1 z2 converges, then lim zm 0 (note that if lim zm 0
m m
then a series z1 z2 may or may not converge). Hence lim zm 0 then the series diverges.
m
Power Series: Power series are the most important series in complex analysis because their sums are
analytic function and are represented by power series. A power series in powers of z z0 is a series
of the form,
n0 an ( z z0 )n a0 a1 ( z z0 ) a2 ( z z0 )2 (4.43)
where z is a complex variable; a0 , a1 , are complex or real constants, called the coefficients of the
series; and z0 is a complex or real constant, called the centre of the series. This generalizes real power
series of calculus. If z0 0 , we obtain as a particular case a power series in powers of z , as
n 0 an z n a0 a1 z a2 z 2 (4.44)
Radius of convergence of power series: We consider the smallest circle with centre z0 that
includes all the points at which a given power series (Eq. 4.43) converges. Let R denote its
radius. The circle z z0 R is called the circle of convergence and its radius R the radius of
convergence of Eq. 4.43. Hence, convergence of power series implies that
Convergence everywhere within that circle, i.e., for all z for which z z0 R (the open
disk with centre z0 and radius R ).
As R is as small as possible, the series (Eq. 4.43) diverges for all z for which z z0 R
On the circle the series (Eq. 4.43) may converge at some or all or none of these points.
We denote R if the series (Eq. 4.43) converges for all z ; and R 0 if the series (given
by Eq. 4.43) converges only at the centre z0 .
Example 4.47 [IN-2016 (1 mark)]: In the neighbourhood of z 1 , the function f ( z ) has a power
series expansion of the form f ( z ) 1 (1 z ) (1 z ) 2 . Then f ( z ) is
(a) 1 z (b) 1 ( z 2) (c) ( z 1) ( z 1) (d) 1 (2 z 1)
1 1 1
Solution (a): As 1 z z 2 z 3 , so 1 (1 z ) (1 z ) 2 f ( z) .
1 z 1 (1 z ) z
Taylor and Maclaurin Series: The Taylor series of a function f ( z ) is given as,
f ( z ) n 1 an ( z z0 ) n (4.45)
1 1 f ( z*)
2 i
where, an f ( n ) ( z 0 ) or an C ( z * z )n 1 dz * (from Eq. 4.40). If we integrate Eq. 4.40
n! 0
counter-clockwise around a simple closed path C that contains z0 in its interior and is such that
f ( z ) is analytic in a domain containing C and every point inside C . A maclaurin series is a Taylor
series with centre z0 0 . The remainder of the Taylor series (Eq. 4.45) after the term an ( z z0 ) n is
( z z0 ) n 1 f ( z*)
Rn ( z )
2 i
C ( z * z0 )n1 ( z * z ) dz * (4.46)
This series converges and represents f ( z ) in the enlarged open annulus obtained from the given
annulus by continuously increasing the outer circle C1 and decreasing C 2 until each of the two circles
reaches a point where f ( z ) is singular. In the important special case that z0 is the only singular point
of f ( z ) inside C 2 , this circle can be shrunk to the point z0 , giving convergence in a disk except at
the centre. In this case the series (or finite sum) of the negative powers of Eq. 4.48 is called the
principal part of the singularity of f ( z ) at z0 .
Instead of Eq. 4.48 and 4.49, we may write, by denoting bn by a n , as
f ( z ) n an ( z z0 ) n (4.50)
where all the coefficients are now given by a single integral formula as,
1 f ( z*)
an
2 i C ( z * z0 ) n 1
dz * , ( n 0, 1, 2, ) (4.51)
Example 4.49 [IN-2007 (2 marks)]: For the function (sin z ) z 3 of a complex variable z , the point
z 0 is
(a) a pole of order 3 (b) a pole of order 2 (c) a pole of order 1 (d) not a singularity
3 5 7 2 4
sin z 1 z z z 1 1 z z
Solution (b): As 3
3 z 2 . As all the negative
z z 3! 5! 7! z 3! 5! 7!
nd
powers of ( z 0) in our expression after 2 are missing, so the singularity at z 0 is called a pole of
order 2.
Here we integrate counter-clockwise around a simple closed path C that contains z z0 in its
interior (but no other singular points of f ( z ) on or inside C ). The coefficient b1 is called the
residue of f ( z ) at z z0 and is denoted by,
b1 Re s f ( z ) (4.54)
z z0
Formulas for Residue: To calculate a residue at a pole, we need not produce a whole Laurent
series, but we can derive formulas for residues once and for all.
Simple Poles: The two formulas for the residue of f ( z ) at a simple pole at z0 are
Res f ( z ) b1 lim ( z z0 ) f ( z ) (4.55)
z z0 z z0
Poles of any order: The residue of f ( z ) at an mth order pole at z0 is given as,
1 d m 1 m
Res f ( z ) lim m1 ( z z0 ) f ( z ) (4.57)
z z0 (m 1)! 0 dz
z z
In particular, for a second – order pole ( m 2) is given as, Res f ( z ) lim
z z0 z z0 d
dz
( z z0 ) 2 f ( z ) .
1 2z 1 2z 1 4 3
Res f ( z ) lim( z 2) f ( z ) lim( z 2) lim
z 2 z2 z ( z 1)( z 2) z 0 z ( z 1)
z2 (2)(1) 2
[Similar question was also asked in CH-2013 (2 marks)]
Example 4.53 [XE-2010 (1 mark)]: The residue of the function (a) 2 (b) 1
f ( z ) (sin 4 z ) ( z 4)3 at z 4 is (c) –1 (d) –2
Solution (b): As the denominator of the given function is zero at z 4 so z 4 is a pole of
1 d2
order 3. Thus from Eq. 4.69, we have Res f ( z ) lim ( z 4)3 f ( z )
2
2! z 4
z 4
dz
1 d 2 3 sin 4 z 1 d2
Res f ( z ) lim 2 ( z 4) 3 lim 2
sin 4 z
z 4 2 z 4
dz z 4 2 z 4 dz
Res f ( z )
z 4
1
2 z
lim
4 d
dz 1
4 sin 3 z cos z lim 4 sin 4 z 12 sin 2 z cos 2 z
2 z 4
Res f ( z )
z 4
1
2
4
1
4
1 1
2 2
1
2
12 1 3 1
sin( z )
Example 4.54 [EC-2016 (1 mark)]: For f ( z ) , the residue of the pole at z 0 is _____.
z2
Solution: As the denominator of the given function is zero at z 0 , which is a pole of order 2. Thus
d 2 1
Res f ( z )
z 0
1
(2 1)! z 0 dz
2 1 d 2 sin( z )
lim 2 1 ( z 0) f ( z ) lim z
1! z 0 dz
lim
z z 0 dz
2
d
sin z
Res f ( z ) lim cos z 1 .
z 0 z 0
Example 4.55 [CE-2005 (2 marks)]: Consider likely applicability of Cauchy’s Integral Theorem to
evaluate the following integral counter clockwise around the unit circle C . I
sec z dz , z being a C
complex variable. The value of I will be
(a) I 0 : singularities set = {(2n 1) 2} , n 0,1, 2, (b) I 0 : singularities set =
(c) I 2 : singularities set = n , n 0,1, 2, (d) None of these
Solution (b):
1 3 3
As I C sec z dz C cos z dz which has poles at z (2n 1) 2 , i.e., z , 2 , 2 , 2 , 2 ,
and none of these poles lie inside z 1 . So we have singularities set = . Also the sum of residues
of the poles which lie inside the given closed region is zero as none of the poles lie inside z 1 .
1
Thus from Eq. 4.58 we have I
sec z dz dz 2 i 0 0 .
C C cos z
1 1
Solution (a): As the given function is f ( s )
. So the pole of the given function
2
s 1 ( s 1)( s 1)
is s 1 , which are simple poles. As the pole s 1 lies outside the given semi-circle; but the pole
s 1 lies within the given semi-circle. So from Eq. 4.58 we have C f ( z)dz 2 i (sum of residues
of the poles which lie inside the given closed region). As we have only one pole s 1 which lie
within the given semi-circle. So C f ( s)ds 2 j Res
z 1
f ( s) . As
( s 1) 1 1 1
Res f ( s ) lim( s 1) f ( s ) lim
z 1 z 1 z 1 ( s 1)( s 1)
lim
z 1 s 1
2
. So C f ( s)ds 2 j 2 j .
Example 4.57 [ME-2008 (2 marks)]: The integral f ( z )dz evaluated (a) 2 i (b) 4 i
around the unit circle on the complex plane for f ( z ) (1 z ) cos z is (c) 2 i (d) 0
Solution: As the given function
cos z 1 z 2 z 4 z6 1 z z 3 z5
f ( z) 1 f ( z ) , which has a simple pole at
z z 2! 4! 6! z 2! 4! 6!
z 0 , which lie inside the unit circle, so from Eq. 4.57, we have
have C f ( z)dz 2 i (sum of residues of the poles which lie inside the given closed region), i.e.,
cos z
C z
dz 2 i 1 2 i
C g ( z ) dz 2 i (1 c0 ) .
cos(2 z )
Solution (c): As the given function is f ( z ) . So the pole of the given function is
(2 z 1)( z 3)
z 1 2 , 3 , which all are simple poles. As the pole z 3 lie outside the circle z 1 since 3 1 ; but
the pole z 1 2 lie inside the circle z 1 since 1 2 1 . Thus from Eq. 4.55, we have
cos(2 z ) cos(2 z ) 1
Res f ( z ) lim ( z 1 2) f ( z ) lim ( z 1 2) zlim ;
z 1 2 z 1 2 z 1 2
(2 z 1)( z 3) 1 2 2( z 3) 5
So from Eq. 4.58 we have C f ( z)dz 2 i (sum of residues of the poles which lie inside the given
cos(2 z ) 1 2 i
closed region). Thus C (2 z 1)( z 3) dz 2 i 5 5
.
Example 4.60 [CH-2009 (2 marks)]: Using the residue theorem, the value of the integral (counter-
clockwise) (8 7 z ) ( z 4) dz around a circle with centre at z 0 and radius 8 (where z is a
z 4 z 4 z4
Res f ( z ) lim( z 4) f ( z ) lim ( z 4)
8 7z
z4 lim(8 7 z ) 20 .
z 4
Thus
Solution: As the given function is f ( z ) ( z 2 8) (0.5 z 1.5i ) . So the pole of the given function is
z 3i , which is a simple poles. As the poles z 3i lies inside the given circle z 4 (since 3i 4 ).
So from Eq. 4.58 we have C f ( z)dz 2 i (sum of residues of the poles which lie inside the given
closed region). As we have only one pole z 3i lie within the circle z 4 . So
3 z 4 3z 4
Solution: As the given function is f ( z ) 2
. So the pole of the
( z 4 z 5) ( z 2 i )( z 2 i )
given function is z 2 i , which all are simple poles; but both the poles are not les within the circle
z 1 , as for both poles we have z 1 . So the sum of the residue of poles which lies inside the
given closed region is zero; and thus from Eq. 4.58 we have C f ( z)dz 2 i 0 0 .
C e
1z
Example 4.63 [IN-2011 (1 mark)]: The contour integral dz with C as the counter-clockwise
unit circle in the z plane is equal to
(a) 0 (b) 2 (c) 2 1 (d)
1 1 1
Solution (c): As f ( z ) e1 z 1 which has isolated essential singularity at
2
z 2! z 3! z 3
z 0 [since we have infinitely many negative terms of ( z 0) ]. So from Eq. 4.54, we have the
residue of f ( z ) is the coefficient of the first negative power of ( z 0) , which is 1. So
C e
1z
dz 2 i Res 2 i 1 2 i 2 1 .
z 0
C z
2
Example 4.64 [PI-2011 (2 marks)]: The value of ( z 4 1) dz using Cauchy’s integral around
the circle z 1 1 , where z x iy is
(a) 2 i (b) i 2 (d) 2 i
(c) 3 i 2
Solution: As the given function is f ( z ) z 2 ( z 4 1) z 2 {( z 1)( z 1)( z i )( z i )} . So the pole of
the given function is z 1, i , which all are simple poles. As the poles z 1, i lies outside the
given circle (since 1 1 1 , i 1 1 , i 1 1 ); but the pole z 1 lies within the circle z 1 1
, as 1 1 1 . So from Eq. 4.58 we have C f ( z)dz 2 i (sum of residues of the poles which lie
inside the given closed region). As we have only one pole z 1 lie within the circle z 1 1 . So
Example 4.65 [CH-2012 (2 marks)]: If i 1 , the value of the integral C [(7 z i) {z ( z 2 1)}]dz
, z 2 using the Cauchy residue theorem is
(a) 2 i (b) 0 (c) 6 (d) 6
7z i 7z i
Solution (b): As the given function is f ( z ) 2
. So the pole of the given
z ( z 1) z ( z i)( z i)
function is z 0, i , which all are simple poles, and all lies within the circle z 2 , since for all
poles we have z 2 . So from Eq. 4.58 we have C f ( z)dz 2 i (sum of residues of the poles
which lie inside the given closed region). Thus from Eq. 4.55, we have
7z i 7z i i
Res f ( z ) lim( z 0) f ( z ) lim ( z 0) 2 lim 2 i;
z 0 z 0 z 0
z ( z 1) z 4
( z 1) 1
7z i 7 z i 6i
Res f ( z ) lim ( z i) f ( z ) lim ( z i ) lim 3i
z i z i z i
z ( z i )( z i) z i
z ( z i) 2
7z i 7 z i 8i
Res f ( z ) lim( z i) f ( z ) lim ( z i) lim 2 4i ;
z i z i z i
z ( z i )( z i) z i
z ( z i ) 2i
So sum of residues of the poles which lie inside the given closed region is i 3i 4i 0 . Thus
C f ( z)dz 2 i (0) 0 .
[Similar questions were also asked in XE-2007 (2 marks)]
1 2
Example 4.66 [EC-2012, EE-2012, IN-2012 (1 mark)]: Given f ( z ) . If C is a
z 1 z3
counter-clockwise path in the z plane such that z 1 1 , the value of {1 (2 i )}
f ( z)dz is C
(a) –2 (b) –1 (c) 1 (d) 2
Solution (c): From Eq. 4.58 we have {1 (2 i)}
f ( z )dz sum of residues of the poles which lie
C
1 2 z 1
inside the given closed region. As the given function is f ( z ) . So the
z 1 z 3 ( z 1)( z 3)
poles of the given function are z 1, 3 and both are simple poles. As the simple pole z 1 lies
within the given region z 1 1 , since 1 1 1 ; also the simple pole z 3 lies outside the given
region z 1 1, since 3 1 1 so
1
f ( z )dz Res
2 i C z 1
f ( z ) lim ( z 1) f ( z ) lim ( z 1)
z 1 z 1
z 1
lim
( z 1)( z 3) z 1
z 1
z 3
2
2
1.
( z
2
Example 4.67 [EE-2013 (2 marks)]: 4) ( z 2 4) dz evaluated anticlockwise around the
circle z i 2 , where i 1 , is
(a) 4 (b) 0 (c) 2 (d) 2 2i
2 2
z 4 z 4
Solution (a): As the given function is f ( z ) 2
. So the pole of the given
z 4 ( z 2i )( z 2i)
function is z 2i , which are simple poles. As the pole z 2i lies outside the given circle (since
2i i 2 ); but the pole z 2i lies within the circle z i 2 , as 2i i 2 . So from Eq. 4.58 we
have C f ( z)dz 2 i (sum of residues of the poles which lie inside the given closed region). As we
have only one pole z 2i lie within the circle z i 2 . So f ( z ) dz 2 i Res f ( z ) . As
C z 2 i
( z 2i )( z 2 4) ( z 2 4) 8
Res f ( z ) lim ( z 2i ) f ( z ) lim Res f ( z ) lim 2i . So
z 2 i z 2i z 2i ( z 2i )( z 2i ) z 2i z2i z 2i 4i
but the pole z 1 lies within the circle z 1 1 , as 1 1 1 . So from Eq. 4.58 we have
C f ( z)dz 2 i (sum of residues of the poles which lie inside the given closed region). As we have
only one pole z 1 lie within the circle z 1 1 . So f ( z ) dz 2 i Res f ( z ) . As
C z 1
( z 1) z 2 z2 1 1
Res f ( z ) lim( z 1) f ( z ) lim
z 1 z 1 z 1 ( z 1)( z 1)
lim
z 1 z 1
2
. So C f ( z )dz 2 i 2 i .
2z 5
Example 4.69 [EE-2016 (1 mark)]: The value of the integral C {z (1 2)}( z 2 4 z 5) dz over the
2z 5 2z 5
f ( z )dz 2 i lim{z (1 2)} 2
2 i lim 2
C z 1 2 {z (1 2)}( z 4 z 5) z 1 2 ( z 4 z 5)
2(1 2) 5 48
f ( z )dz 2 i 2
i
C (1 2) 4(1 2) 5 13
Example 4.70 [EC-2016 (2 marks)]: In the following integral, the contour C encloses the points
1 sin z
2 j and 2 j . The value of the integral
2 C ( z 2 j )
3
dz is _____.
sin z
Solution: For the given function f ( z ) 3
, the denominator is zero at z 2 j , so z 2 j
( z 2 j )
1 sin z 1
is a pole of order 3. So I
2
C ( z 2 j )
3
dz
2
2 j (sum of residues of the poles which
1 ez
2 j
C z 2 dz along a closed
Example 4.71: [EC-2016 (2 marks)]: The values of the integral
contour C in anti-clockwise direction for (i) the point z0 2 inside the contour C , and (ii) the point
z0 2 outside the contour C , respectively, are
(a) (i) 2.72, (ii) 0 (b) (i) 7.39, (ii) 0 (c) (i) 0, (ii) 2.73 (d) (i) 0, (ii) 7.39
z
e
Solution (b): As the given function f ( z ) has the pole where z 2 0 z 2 is a pole,
z2
which is of order 1 and lies inside the closed contour C ; so
(i) For the point z0 2 inside the contour C
1
2 j
C f ( z)dz (sum of residues of the poles which lies inside the given closed region).
1 ez
2 j
C f ( z ) dz Res f ( z ) lim( z 2) lim e z e 2 7.389 .
z 2 z2 z2 z 2
1
(ii) For the point z0 2 outside the contour C , so Res f ( z ) 0 and
z 2 2 j
C f ( z)dz (sum of
1
residues of the poles which lies inside the given closed region)
2 j
C f ( z)dz 0 .
1 z2 1
2 j C z 2 1
Example 4.72 [IN-2016 (2 marks)]: The value of the integral dz , where z is a
complex number and C is a unit circle with centre at 1 0 j in the complex plane is _____.
Solution: For f ( z ) ( z 2 1) ( z 2 1) , pole of the given function are the points where
z 2 1 z 1,1 , which are all simple poles. As C is a unit circle with centre at 1 0 j in the
complex plane, so for z 1 , 1 (1 0 j ) 2 2 1 z 1 lies outside the contour C .
for z 1 , 1 (1 0 j ) 0 1 z 1 lies inside the contour C .
1 z2 1
2 j C z 2 1
Hence dz (sum of residues of the poles which lies inside the contour C ).
1 z2 1 z2 1 z2 1 2
2 j C z 2 1
dz Res f ( z ) lim{z 1} f ( z ) lim{z 1} 2
lim 1
z 1 z 1 z 1 z 1 z 1 z 1 2
sin x
Example 4.73 [ME-2016 (2 marks)]: The value of the integral x 2 2 x 2 dx evaluated using
contour integration and the residue theorem is
(a) sin(1) e (b) cos(1) e (c) sin(1) e (d) cos(1) e
iz
sin z Im(e )
Solution (a): Let f ( z) 2
2
, since eiz cos z i sin z Im(eiz ) sin z ,
z 2z 2 z 2z 2
2
whose poles are values of ‘ z ’ such that z 2 z 2 0 z 1 i , both of which are simple poles;
but only z 1 i lie in upper half of Z-plane.
f ( z ) dz (2 i) (sum of residues of the poles which lies inside the given closed region). As
f ( z ) dz (2 i) sin(1)
2ie
sin(1)
e
. Thus option (a) is correct.
3z 5
Example 4.74 [ME-2016 (2 marks)]: The value of ( z 1)( z 2) dz along a closed path to
Solution (b): f ( z)dz (2 i) (sum of residues of the poles which lies inside given closed region)
3z 5
As ( z 1)( z 2) dz 4 i 2 i (2) sum of residues of the poles which lies inside given
1 1
1 1
5
, which is in the form of
1 1
which is 1 (4!) 0.0416 . Hence I
2 i
z 7 cos 2 dz 0.0416 .
z
Exercise: 4.2
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. Evaluate the complex integral C zdz , where C is the straight line path from z 1 i to z 3 i
(a) 4 2i (b) 4 2i (c) 2 i (d) 2 i
2. Evaluate the complex integral C zdz , where C is the straight line path from z 3 i to
z 3 3i .
(a) 4 6i (b) 4 6i (c) 4 6i (d) 4 6i
3. Evaluate the complex integral C zdz , where C is the straight line path from z 1 i to
z 3 3i .
(a) 8 (b) 8 (c) 8i (d) 8i
2
4. Evaluate the line integral C z dz , where C is a part of unit circle in anticlockwise direction from
z 1 to z i .
(a) (1 i) 3 (b) (1 i ) 3 (c) (1 i) 3 (d) (1 i ) 3
5. Evaluate the line integral C (1 z)dz , where C is a unit circle.
(a) 2 i (b) 2 i (c) i (d) i
2
6. Evaluate the line integral C ( z z ) dz , where C is a part of unit circle in anticlockwise direction
from z 1 to z i .
(a) (4 3) i(1 3) (b) (4 3) i (1 3) (c) (4 3) i (1 3) (d) (4 3) i(1 3)
7. Evaluate the line integral C {1 ( z z0 )}dz , where C is a circle centred at z0 and of any radius.
The path traced out once in the anticlockwise direction.
(a) 2 i (b) 2 i (c) i (d) i
2
8. Evaluate the line integral C z dz , where the curve C is the line segment with initial point 1
and final point i .
(a) (3 2)(1 i) (b) (2 3)(1 i ) (c) (3 2)(1 i ) (d) (2 3)(1 i)
2
9. Evaluate the line integral C z dz , where the curve C is the arc of the unit circle Im( z ) 0 with
initial point 1 and final point i .
(a) 1 i (b) 1 i (c) 1 i (d) 1 i
2 i
10. Suppose is real, then e 1 is
(a) 2 (b) 2 (c) 2 (d) 2
2
11. If C is the line segment joining ( 1 i ) and (1 i) , then C (1 z ) is
(a) 2 (b) 2 (c) 2 (d) 2
z 2
12. The maximum value of C {e ( z 1)} , where C is the circle z 2 traversed once in the
counter-clockwise direction, is
(a) (4 3) e (b) (3 4) e (c) (3 4) e 2 (d) (4 3) e 2
13. If C is the curve y x3 3 x 2 4 x 1 joining points (1,1) and (2,3) , the value of
2
C (12 z 4iz ) dz is
(a) 156 38i (b) 156 38i (c) 156 38i (d) 156 38i
14. Evaluate f ( z )
{e 2z
( z 1) }dz 4
z 3
C (1 z
6
20. Evaluate ) sin zdz , where C is a circle of unit radius centred at origin.
(a) i 60 (b) i 120 (c) i 30 (d) i 90
2 1 1 3
21. The singularities of the function f ( z ) 4
are
2
z z z i ( z i)
(a) A pole of order 2 at z 0 ; a simple pole at z i ; and a pole of order 4 at z i
(b) A simple at z 0, i ; and a pole of order 4 at z i
(c) A simple pole at z 0, i ; and a pole of order 2 at z i
(d) A pole of order 2 at z 0 ; a simple pole at z i ; and a pole of order 4 at z i
1 1 2
22. The singularities of the function f ( z ) 2 2
3
are
z ( z i) ( z i)
(a) A double pole at z 0 and a pole of order 2 at z i
(b) A pole at z 0 and a pole of order 3 at z i
(c) A double pole at z 0 and a pole of order 3 at z i
(d) A pole at z 0 and a pole of order 2 at z i
23. The residue of f ( z ) 1 ( z 2 4) at z 2i and z 2i are, respectively,
(a) 1 (4i) , 1 (4i ) (b) 1 (4i ) , 1 (4i) (c) 4i , 4i (d) 4i , 4i
2 2
24. The singularity of f ( z ) 1 {z ( z 9)} is
(a) A pole of order 4 at z 0 ; Simple pole at z 3i
(b) Double pole at z 0 ; Simple pole at z 3i
(c) Double pole at z 3i ; Simple pole at z 0
(d) A pole of order 4 at z 3i ; Simple pole at z 0
25. The residue at each of the singularity of the function f ( z ) 1 {z 2 ( z 2 9)} is
(a) Res z 3i i 54 , Res z 3i i 54 and Res z 0 i
(b) Res z 3i i 54 , Res z 3i i 54 and Res z0 0
(c) Res z 3i i 54 , Res z 3i i 54 and Res z0 0
(d) Res z 3i i 54 , Res z 3i i 54 and Res z 0 i
C {1 ( z
2
26. Evaluate the integral 4)}dz , where C is the circle z 2i 1
(a) 2 (b) 2 (c) 2 (d) 0
C {1 ( z
2
27. If C is any closed path enclosing both z 2i and z 2i , then 4)}dz _____.
5z 2
28. z 3 z ( z 1) dz k i , where k _____.
C e
1 z
29. dz k i , where k _____.
30. The solutions of z 3 i is
i ( 6) i ( 3) i (3 i ( 6)
(a) z e ,e ,e 2)
(b) z e , ei (5 6) , ei ( 2)
8
33. All the solutions of the equation z 1 is
1 1 1 1 1 1 1 1
(a) z 1, 2i, i, i (b) z 1, i , i , i
2 2 2 2 2 2 2 2
1 1 1 1 1 1 1 1
(c) z 1, i, i, i (d) z 2, 2i, i, i
2 2 2 2 2 2 2 2
34. Which of the following functions are analytic?
5 iz 2
(i) f ( z ) ze ; (ii) f ( z ) x 2 y ixy 2 ; (iii) f ( z ) x 2 y 2 i (2 xy )
(a) Only (ii) and (iii) (b) Only (i) and (ii) (c) Only (i) and (iii) (d) All (i), (ii) and (iii)
35. Find an analytic function f ( z ) , if such function exists, so that the imaginary part of f ( z ) is
equal to v 3 x 2 4 y 3 y 2 .
(a) f ( z ) (6 xy 4 x c ) i (3 x 2 4 y 3 y 2 ) (b) f ( z ) (6 xy 4 x c) i (3 x 2 4 y 3 y 2 )
(c) f ( z ) (6 xy 4 x c ) i (3 x 2 4 y 3 y 2 ) (d) f ( z ) (6 xy 4 x c ) i(3x 2 4 y 3 y 2 )
where c is any constant.
36. For the function f ( z ) (1 cos z ) z 2 , then z 0 is
(a) a removable singularity (b) an essential singularity
(c) an isolated singularity (d) a pole of order 2
37. For the function f ( z ) z cos(1 z ) , then z 0 is
(a) a removable singularity (b) an essential singularity
(c) an isolated singularity (d) a pole of order 2
38. Evaluate the complex integral of given function f ( z ) z cos(1 z ) over the square with vertices
(1, 0) , (0,1) , ( 1, 0) and (0, 1) .
(a) i (b) i (c) 2 i (d) 2 i
Answer Keys
Answer Keys: Exercise: 4.1
1 2 3 4 5 6 7 8 9 10 11 12 13 14 15
0 a 0 d c d d 1 a d d a a a a
16 17 18 19 20 21 22
b b c b b 3 c
[AG – 2017 (2 marks)]: The areas of seven horizontal cross-sections of a water reservoir at intervals
of 9 m are 210, 250, 320, 350, 290, 230 and 170 m2. The estimated volume of the reservoir in m3
using Simpson’s rule is _____.
Solution: Applying Simpson’s 1/3 rule,
the required volume (9 3){210 4(250 350 230) 2(320 290) 170} 14760 m3.
2
For Exact Solution: On integrating the given DE, we get: du (3t 1)dt u t 3 t c , on
applying: u 0 at t 0 , we get c 0 . So exact solution is given by: u (t ) t 3 t . So
u (2) exact 23 2 10 .
Hence absolute error is given by: u (2)exact u (2) approx. 10 2 8 .
[CH – 2017 (2 marks)]: Match the problem type in Group-I with the numerical method in Group-2.
Group-I Group-II Choose the correct set of
P) System of linear algebraic equations I) Newton-Raphson combinations:
Q) Non-linear algebraic equations II) Gauss-Seidel (a) P-II, Q-I, R-III, S-IV
R) Ordinary differential equations III) Simpson’s rule (b) P-I, Q-II, R-IV, S-III
S) Numerical integration IV) Runge-Kutta (c) P-IV, Q-III, R-II, S-I
(d) P-II, Q-I, R-IV, S-III
Solution (d): Gauss-Seidel iterative method is used for solving a set of linear algebraic equation only.
Newton-Raphson method is used for solving non-linear algebraic equation. Among the given options
in Group-II, Runge-Kutta method is used for solving ODEs; and Numerical integration problems are
solved by using Simpson’s rule.
[EC – 2017 (2 marks)]: Starting with x 1 , the solution of the equation x 3 x 1 , after two
iterations of Newton-Raphson’s method (up to two decimal places) is _____.
[EE – 2017 (2 marks)]: Only one of the real roots of f ( x ) x 6 x 1 lies in the interval 1 x 2
and bisection method is used to find its value. For achieving an accuracy of 0.001, the required
minimum number of iterations is _____.
ba
Solution: We know that for bisection method , where a and b are lower and upper limit
2n
of the interval, so a 1 , b 2 ; n is the number of iterations required to achieve an accuracy of
2 1
0.001 . Thus n
0.001 2 n 103 n log10 2 3 n 3 log10 2 n 9.965 . So
2
minimum number of iterations is n 10 .
[MA – 2017 (2 marks)]: Let p( x ) be the polynomial of degree at most 2 that interpolates the data
( 1, 2) , (0,1) and (1, 2) . If q ( x ) is a polynomial of degree at most 3 such that p ( x ) q ( x )
interpolates the data ( 1, 2) , (0,1) , (1, 2) and (2,11) , then q(3) equals _____.
Solution: p (3) L0 (3) f0 L1 (3) f1 L2 (3) f 2 , where
{3 0}{3 1} {3 ( 1)}{3 1} {3 ( 1)}{3 0}
L0 (3) 3 , L1 (3) 8 ; L2 (3) 6
{1 0}{1 1} {0 ( 1)}{0 1} {1 ( 1)}{1 0}
So p (3) 3 2 ( 8) 1 6 2 6 8 12 10 .
[ME – 2017 (2 marks)]: P (0,3) , Q (0.5, 4) and R (1,5) are three points on the curve defined by
f ( x ) . Numerical integration is carried out using both Trapezoidal rule and Simpson’s rule within
limits x 0 and x 1 for the curve. The difference between the two results will be
(a) 0 (b) 0.25 (c) 0.5 (d) 1
Solution: Using Trapezoidal rule: ITrapezoidal 0.5 (1 2) f (0) f (0.5) (1 2) f (1)
ITrapezodal 0.5 (1 2) 3 4 (1 2) 5 4 .
Using Simpson’1 1/3 rule: I Simpson 1 3 Rule (0.5 3) f (0) 4 f (0.5) f (1)
I Simpson 1 3 Rule (0.5 3) 3 4 4 5 4 . So ITrapezodal I Simpson 1 3 Rule 4 4 0 .
1 x2
[MT – 2017 (2 marks)]: The definite integral, 0 e dx is to be evaluated numerically. Divide the
integration interval into exactly 2 subintervals of equal length. Applying the trapezoidal rule, the
approximate value of the integral is _____ (answer up to two decimal places).
x2
Solution: We have f ( x ) e , and h (1 0) 2 0.5 .
1 x2
So Using Trapezoidal rule: 0 e 0.5 (1 2) f (0) f (0.5) (1 2) f (1)
1 2 2 2 2
e x 0.5{(1 2)e 0 e 0.5 (1 2)e 1 } 0.5{0.5 0.78 0.18} 0.73 .
0
[MT – 2017 (2 marks)]: Using the bisection method, the root of the equation: x 3 x 1 0 , after
three iteration is _____ (answer up to two decimal places; assume starting values of x 1 and 1 ).
Solution: Let f ( x ) x 3 x 1 , as f ( 1) 3 0 and f (1) 1 0 , so there will be a root lie in
( 1,1) .
First iteration: x0 (1 1) 2 0 . As f (0) 1 0 and f (1) 1 0 , so there will be a root lie in
(0,1)
Second iteration: x1 (0 1) 2 0.5 . As f (0.5) 0.375 0 and f (1) 1 0 , so there will be a
root lie in (0.5,1) .
Third iteration: x2 (0.5 1) 2 0.75 . As f (0.75) 0.172 0 and f (0.5) 0.375 0 , so there
will be a root lie in (0.5, 0.75) .
Thus after 3rd iteration, the root of the equation is x2 0.75 .
d2y
[PE – 2017 (1 mark)]: If f ( x, y ) 0 is to be solved using the conditions y (0) a and
dx 2
y (1) b , which of the following numerical method(s) can be used?
(a) Euler with shooting method (b) Euler without shooting method
th
(c) 4 order Runge-Kutta with shooting method (d) Both (a) and (c)
Solution (d): Euler and Runge-Kutta method are used for solving differential equations. But the
numerical solution of differential equation with boundary conditions are solved by Shooting method.
So options (a) and (c) are correct.
[PE – 2017 (1 mark)]: The numerical method used to find the root of a non-linear algebraic equation,
that converges quadratically is:
(a) Bisection method (b) Regula-falsi method (Method of False Position)
(c) Newton-Raphson method (d) None of the above
Solution (c): Newton – Raphson process has a second order or quadratic convergence
1
[PE – 2017 (2 marks)]: Single step integration (step size 0.5) of I x 2 e x dx , evaluated
0
numerically using the Simpson’s 1 3 rule, is _____. (write answer with three decimal places)
Solution: We have f ( x ) x 2 e x , and h 0.5 ; so using Simpson’s 1 3 rule:
1
I x 2 e x (h 3){ f (0) 4 f (0.5) f (1)} (0.5 3){0 4 0.52 e 0.5 1 e1} 0.728 .
0
[PE – 2017 (2 marks)]: Solve dy dx y numerically from x 0 to 1 using explicit, forward, first
order Euler method with initial condition of y (0) 1 and step size ( h) of 0.2. The absolute value of
error in y (1) calculated using analytical and numerical solution is _____ % (calculate the error using
analytical solution as the basis and use three decimal places).
dy
Solution: Analytic solution: dy dx y dx ln y x c ; Applying y (0) 1 , we
y
get c 0 . So ln y x y ( x ) e x . Thus y (1) e 1 0.368
Numerical Solution: yn 1 yn h (dy dx ) ( x ; where xn x0 nh
n , yn )
[PI – 2017 (1 mark)]: Using Simpson’s 1/3 rule for numerical integration, the consecutive points are
joined by a
(a) Line (b) parabola (c) polynomial with power 3 (d) polynomial with power 1 3
Solution (b): In Simpson’s 1/3 rule for numerical integration, the consecutive points are joined by a
parabola.
[PI – 2017 (2 marks)]: Runge-Kutta fourth order method is used to solve the differential equation
dy dx y x . If the initial value y (0) 2 and step-size is 0.1, then the value of y (0.1) is _____ (up
to three decimal places).
Solution: dy dx f ( x, y) y x ; and h 0.1 , so from Eq. 5.60, we have
y1 y0 (1 6)( k1 2k2 2k3 k4 ) , where x1 x0 h ,
k1 hf ( x0 , y0 ) h( y0 x0 ) 0.1(2 0) 0.2 ;
k2 hf {x0 (1 2) h, y0 (1 2)k1} 0.1 f {0.05, 2.1} 0.1 (2.1 0.05) 0.205
k3 hf {x0 (1 2) h, y0 (1 2) k2 } 0.1 f {0.05, 2.1025} 0.1 (2.1025 0.05) 0.20525
k3 hf {x0 h, y0 k3 } 0.1 f {0.1, 2.20525} 0.1 (2.20525 0.1) 0.210525
[TF – 2017 (2 marks)]: Using Simpson’s 1/3 rule, the value of the integral (1 ) (1 sin x ) dx ,
0
accurate to two decimal places, is _____.
Solution: Let f ( x ) 1 sin x ; and h ( 0) 2 2 . So
0 (1 sin x ) dx {( 2) 3}{ f (0) 4 f ( 2) f ( )} ( 6){1 4 2 1} (5 3)
Thus (1 ) (1 sin x ) dx 5 3 1.66 .
0
Numerical methods are used to solve problems on computers or calculators by numeric calculations,
resulting in a table and/or graphical representation (figures). Before going to discuss the numerical
methods, let us discuss the following terms:
Significant Digits: The digits that are known reliably plus the first uncertain digit are known as
significant digits or figures. For e.g., when a measured distance is reported to be 364.7 m, it has four
significant digits 3, 6, 4 and 7; the digits 3, 6 and 4 are certain and reliable, but the digit 7 is uncertain.
Following are some of the common rules for counting significant digits in a reported measurement:
1. All non-zero digits are significant. For e.g., x 1236 and x 1.236 both have 4 significant
digits.
2. All zeros occurring between two non-zero digits are significant digits, no matter where the
decimal point is in the number. For e.g., x 1007 and x 1.0809 have 4 and 5 significant digits,
respectively.
3. If the number is less than one, the zero(s) on the right of the decimal point and to left of first non-
zero digit are not significant. For e.g., x 0.005704 have 4 significant digits.
4. In a number without a decimal point, the terminal or trailing zeros are not significant. For e.g.,
x 3210 have 3 significant digits.
5. In a number with decimal point, the trailing zeros are significant digits. For e.g., x 3.210 has 4
significant digits.
6. Change of units does not change the number of significant figures in a measurement. For e.g., the
length x 5.608cm 56.08mm cm has 4 significant digits
7. The digit 0 which is put on the left of a decimal for a number less than 1 is never significant. For
e.g., 0.3570 has 4 significant digits.
Note: Referring to rules 4, 5 and 6, there can be some confusion regarding the trailing zeros. Suppose
a measured length is reported as x 4.700 m, which has four significant figures. On changing the
units, we can rewrite the same length as x 0.004700 km or x 470.0 cm or x 4700 mm. As per
rule 4, we conclude that x 4700 m has 2 significant digits. The fact that x 4700 mm has 4
significant digits, as change of unit cannot change the number of significant digits. To remove such
ambiguities in determining the number of significant digits, the best way is to report every
measurement in scientific notation, i.e., in power of 10, often putting the decimal after the first digit.
For e.g. x 4.700m 4.700 10 2 cm 4.700 10 3 mm 4.700 10 3 km , in each case, the number of
significant figures is 4, as the power of 10 is irrelevant to the determination of significant figures;
hence the confusion disappears.
Rounding Off of Numbers: The result of computation with approximate numbers, which contain
more than one uncertain digit should be rounded off. If a number is to be rounded off to n significant
digits, then we follow the following rules:
1. Discard all digits to the right of the nth digit if ( n 1) th digit is less than 5. For e.g., x 7.82 is
rounded off to two significant digits as 7.8.
2. If the ( n 1) th digit is greater than 5 or it is 5 followed by a nonzero digit, then nth digit is
increased by 1. For e.g., x 6.87 is rounded off to two significant digits as 6.9. If the ( n 1) th
digit is less than 5, then digit remains unchanged. For e.g., x 6.84 is rounded off to two
significant digits as 6.9.
3. If the ( n 1) th digit is 5 and is followed by zero or zeros, then nth digit is increased by 1 if it is
odd and it remains unchanged if it is even. For e.g., x 3.250 is rounded off to two significant
digits as 3.2; x 3.350 is rounded off to two significant digits as 3.4.
Error due to rounding off numbers: If a number is rounded off according to the rules, the
maximum error due to rounding does not exceed the one half of the place value of the last retained
digit in the number. The difference between a numerical value X and its rounded value X 1 is called
round off error is given by E X X 1 .
Truncation and Error due to Truncation of Numbers: Leaving out the extra digits that are
not required in a number without rounding off, is called truncation or chopping off. The difference
between a numerical value X and its truncated value X 1 is called truncation error and is given by
E X X 1 . The maximum error due to truncation of a number cannot exceed the place value of the
last retained digit in the number.
In truncation the numerical value of a ve number is Approximated number
decreased and that of a negative number is increased. obtained by
Number
If we round off a large number of ve numbers to the Chopping Rounding
same number of decimal places, then the average off off
error due to rounding off is zero. 0.335217... 0.3352 0.3352
In case of truncation of a large number of ve 0.666666... 0.6666 0.6667
numbers to the same number of decimal places the
0.123451... 0.1234 0.1235
average truncation error is one half of the place value
of the last retained digit. 0.213450... 0.2134 0.2134
If the number is rounded off and truncated to the 0.213950... 0.2139 0.214
same number of decimal places, then truncation error 0.335750... 0.3357 0.3358
is greater than the round off error.
0.999999... 0.9999 1
Round of error may be ve or ve but truncation
error is always ve in case of ve numbers and ve 0.555555... 0.5555 0.5556
in case of ve numbers.
Relative and Percentage Error of Numbers: The difference between the exact value of a number
X and its approximate value X 1 , obtained by rounding off or truncation, is known as absolute error.
The quantity X X1 X is called the relative error and is denoted by ER . Thus
ER X X 1 X X X . This is a dimensionless quantity. The quantity X X 100 is
known as percentage error and is denoted by E p , i.e. E p X X 100 .
If a number is rounded off to n decimal digits, then | E R | 0.5 10 n 1
If a number is truncated to n decimal places, then | E R | 10 n 1
Example 5.1 [EC-2013 (1 mark)]: The maximum value of until which the approximation
sin holds to within 10% error is
(a) 10o (b) 18o (c) 50o (d) 90o
Solution (b): 10o 0.1745c ; 18o 0.3142 c ; 50o 0.8727 c ; 90o 1.5708c . As
sin is the exact value and is the approximate value. So percentage error, E p sin sin .
o
So, for option (a), E p sin10 0.1745 sin10
o
100 0.4905% 10% .
o
For option (b), E p sin18 0.3142 sin18
o
100 1.6773% 10% .
For option (c), E p sin 50
o
0.8727 sin 50 100 13.9229% 10% .
o
Thus maximum value of until which the approximation sin holds to within 10% error is 18o .
3 2
b b b b
x y to obtain, a y b y c y d 0 .
3a 3a 3a 3a
3 b2 2b 3 bc 3 1 b2 1 2b 3 bc
ay c y d 2
0 y c y d 2
0
3a 27a 3a a 3a a 27 a 3a
1 b2 1 2b3 bc
y 3 ey f 0 , where e c and f d 2
.
a 3a a 27a 3a
3
Now substituting y z ( s z ) in y 3 ey f 0 z ( s z ) e z ( s z ) f 0
z 6 (3s e) z 4 f z 3 s (3s e) z 2 s 3 0 . Now substituting s e 3 we get,
z 6 f z 3 (e 3 27) 0 w 2 fw (e3 27) 0 (By substituting w z 3 ). Now we have a
quadratic equation and once we obtain the solution to this quadratic equation, back substituting using
the previous substitutions ( w z y x ) to obtain the roots to the general cubic equation. Eq.
w 2 f w (e3 27) 0 gives two roots for w ; and using w z 3 we have three roots for each of the
two roots of w , so we have six values of z which gives six values of y ; but three values of y will
be identical to the other three values. So we get only three values of y and hence three values of x .
Example 5.2: Find a positive root of the equation xe x 1 , which lies between 0 and 1.
Solution: Let f ( x ) xe x 1 . Since f (0) 1 and f (1) 1.718 , it follows that a root lies between 0
and 1. Thus, x0 0.5 . Since f (0.5) is negative, it follows that the root lies between 0.5 and 1. Hence
the new root is 0.75, i.e., x1 0.75 . Using the values of x0 and x1 , we calculate 1 , as
1 ( x1 x0 ) x1 100 33.33% . Proceeding in this way, the following table is constructed where
only the sign of the function value is indicated. The prescribed tolerance is 0.05%.
n a b x f ( x) r (%)
1 0 1 0.5 – 0.1756 –
2 0.5 1 0.75 0.5877 33.33
3 0.5 0.75 0.625 0.1676 20.00
4 0.5 0.625 0.5625 0.0129 11.11
5 0.5625 0.625 0.5938 – 0.0753 5.263
6 0.5625 0.5938 0.5781 0.0305 2.707
7 0.5625 0.5781 0.5703 0.0087 1.368
8 0.5625 0.5703 0.5664 – 0.0020 0.688
9 0.5664 0.5703 0.5684 0.0035 0.352
10 0.5664 0.5684 0.5674 0.0007 0.176
11 0.5664 0.5674 0.5669 – 0.0006 0.088
12 0.5669 0.5674 0.5671 – 0.0001 0.035
Thus after 12 iterations, r , finally satisfies the prescribed tolerance. Hence the required root is 0.567
and it is easily seen that this value is correct to three decimal places.
Example 5.3 [CS-2012 (2 marks)]: The bisection method is applied to compute a zero of the
function f ( x ) x 4 x 3 x 2 4 in the interval [1, 9] . The method converges to a solution after ……
iterations.
Solution: As f (1) 14 13 12 4 1 0 and f (9) 9 4 93 9 2 4 0 so using bisection
method we have a solution in [1, 9] and thus x0 (1 9) 2 5 . Now f (5) 54 53 52 4 0 so
we have a solution in [1, 5] and thus x1 (1 5) 2 3 . Now f (3) 34 33 32 4 0 , so we have a
The Method of False Position: This is the oldest method for finding the real root of a non – linear
equation f ( x ) 0 and closely resembles the bisection method. In this method, also known as regula
falsi or the method of chords, we choose two points a and b such that f (a ) and f (b) are of
opposite signs. Hence, a root must lie in between these points. Now, the equation of the chord joining
the two points [ a, f ( a )] and [b, f (b)] is given by,
y f (a ) f (b ) f ( a )
(5.3)
xa ba
The method consists in replacing the part of the curve between the points [ a, f ( a )] and [b, f (b)] by
means of the chord joining these points, and taking the point of intersection of the chord with the x
axis as an approximation to the root. The point of intersection in the present case is obtained by
putting y 0 in Eq. 5.3. Thus we obtain,
f ( a) a f ( b) b f ( a )
x1 a (b a ) (5.4)
f (b ) f ( a ) f (b) f ( a)
which is the first approximation to the root of f ( x ) 0 . If now f ( x1 ) and f (a ) are of opposite
signs, then the root lies between a and x1 , and we replace b by x1 in Eq. 5.4 and generate the next
approximation. The procedure is repeated till the root is obtained to the desired accuracy. Fig. 5.3
gives a graphical representation of the method. The error criterion
Eq. 5.2 can be used in this case also.
The method may give a false root or may not converge if
either a and b are not sufficiently close to each other or
f ( x ) is discontinuous on [ a, b] .
Geometrically speaking, in this method, part of the curve
between the points A a, f ( a) and B b, f (b) is replaced
by the secant AB and the point of intersection of this secant
with x axis gives an approximate value of the root.
It converges more rapidly than bisection. Figure 5.3: Method of False Position
Example 5.4: Find a real root of the equation f ( x ) x 3 2 x 5 0 by using false position method.
Solution: We find f (2) 1 and f (3) 16 . Hence a 2 , b 3 , and a root lies between 2 and 3.
Eq. 5.4 gives x1 {2(16) 3( 1)} {16 ( 1)} 35 17 2.0588 . Now f ( x1 ) 0.3908 and hence the
root lies between 2.0588 and 3.0. Using Eq. 5.4, we get,
x2 {2.0588(16) 3( 0.3908)} 16.3908 2.0813 . Since f ( x2 ) 0.1472 , it follows that the root
lies between 2.0813 and 3.0. Proceeding in this way, we obtain successively:
x3 2.0896, x4 2.0937, x5 2.0939, x6 2.0943, x7 2.0945, . The correct value is 2.0945…, so
that x7 is correct to five significant figures.
Fixed-Point Iteration for Solving Equation ( ) = : We have so far discussed root – finding
methods, which require the interval in which the root lies. We now describe methods which require
one or more starting values of x . These values need not necessarily bracket the root. The first
iteration method, which requires one starting value of x . To describe this method for finding the roots
of the equation
f ( x) 0 (5.5)
we re-write this equation in the form as,
x g ( x) (5.6)
(certainly g ( x) is non-constant). There are many ways of doing this. For e.g., the equation
3 2 3 12 2 13
x x 1 0 can be expressed as either of the forms: x ( x 1) , x ( x 1) , …. Let x0 be an
approximate value of the desired root . Substituting it for x on the right hand of Eq. 5.6, we obtain
the first approximation as, x1 g ( x0 ) then x2 g ( x1 ) , and in general,
xn 1 g ( xn ) , n 0,1, 2, (5.7)
Some questions are now arise:
1. Does the sequence of approximations x0 , x1 , , xn always converge to some number ?
2. If it does, will be a root of the equation x g ( x ) ?
3. How should we choose in order that the sequence x0 , x1 , , xn converges to the root?
The answer to the first question is negative. As an example we consider the equation x 2 x 1 . If we
take x0 0, x1 2, x2 5, x3 33, , etc. and as n increases xn , increases without limit. Hence the
sequence x0 , x1 , , xn does not always converge and later in this section we state the conditions
which are sufficient for the convergence of the sequence.
The second question is easy to answer, for consider the equation xn 1 g ( xn ) which gives the relation
between the approximations at the nth and ( n 1)th stage. As n increases, the left side tend to the
root , and if g is continuous the right hand side tend to g ( ) . Hence, in the limit, we have
g ( ) which shows that is a root of the equation x g ( x ) .
The answer to the third question is given as: Let x be a root of f ( x ) 0 and let I be an interval
containing the point x . Let g ( x ) and g ( x ) be continuous in I , where g ( x ) is defined by the
equation x g ( x ) which is equivalent to f ( x ) 0 . Then if g ( x ) 1 for all x in I , the sequence of
approximations x0 , x1 , x2 , , xn defined by xn 1 g ( xn ) converges to the root , provided that the
initial approximation x0 is chosen in I .
Example 5.6 [MA-2016 (2 marks)]: For the fixed point iteration xk 1 g ( xk ) , k 0,1, 2, ,
consider the following statements P and Q :
( P ) : If g ( x) 1 (2 x) then the fixed point iteration converges to 2 for all x0 [1,100] .
(Q ) : If g ( x ) 2 x then the fixed point iteration converges to 2 for all x0 [0,100] .
Which of the above statements hold TRUE?
(a) both P and Q (b) only P (c) only Q (d) neither P nor Q
Solution: For fixed point iteration, we know that if x be a root of f ( x ) 0 and let I be an
interval containing the point x . Let g ( x ) and g ( x) be continuous in I , where g ( x ) is defined
by the equation x g ( x) which is equivalent to f ( x ) 0 . Then if g ( x ) 1 for all x in I , the
sequence of approximations x0 , x1 , x2 , , xn defined by xn 1 g ( xn ) converges to the root ,
provided that the initial approximation x0 is chosen in I .
In the given question, for both the statements g ( x ) 1 for all x belongs to given interval. So
both the statements are true.
Newton-Raphson Method: This method is generally used to improve the result obtained by one of
the previous methods. Let x0 be an approximate root of the equation f ( x ) 0 and let x1 x0 h be
the exact root, so that f ( x1 ) 0 . Expanding f ( x0 h) by Taylor's series as,
f ( x0 ) hf ( x0 ) ( h 2 2!) f ( x0 ) 0 . Since h is small, neglecting h 2 and higher powers of h , we
get f ( x0 ) hf ( x0 ) 0 or h f ( x0 ) f ( x0 ) . A closer approximation to the root is given by
x1 x0 f ( x0 ) f ( x0 ) . Similarly, starting with x1 , a still better approximation x2 is given by
x2 x1 f ( x1 ) f ( x1 ) .
In general, xn 1 xn f ( xn ) f ( xn ) (5.8)
which is known as the Newton-Raphson formula or Newton's iteration formula.
Convergence of Newton – Raphson method: If we compare Eq. 5.8 with the relation
xn 1 g ( xn ) of iterative method, we get, g ( x ) x f ( x ) f ( x ) which gives,
f ( x ) f ( x )
g ( x ) (5.9)
f ( x) 2
To examine the convergence we assume that f ( x ) , f ( x ) and f ( x ) are continuous and
bounded on any interval containing the root x of the equation f ( x ) 0 . If is a simple
root, then f ( x ) for some 0 in a suitable neighbourhood of . Within this
2
neighbourhood we can select an interval such that f ( x) f ( x) and this is possible since
f ( ) 0 and since f ( x ) is continuously twice differentiable. Hence, in this interval we have,
f ( x ) 1 (5.10)
Therefore, the Newton – Raphson formula (Eq. 5.8) converges, provided that the initial
approximation x0 is chosen sufficiently close to . When is a multiple root, the Newton –
Raphson method still converges but slowly. To obtain the rate of convergence of the method, we
note that f ( ) 0 so that Taylor’s expansion gives
f ( xn ) ( xn ) f ( xn ) (1 2)( xn ) 2 f ( xn ) 0 from which we get,
f ( xn ) 1 f ( xn )
( xn ) ( xn ) 2 (5.11)
f ( xn ) 2 f ( xn )
From Eq. 5.8 and Eq. 5.11, xn 1 (1 2)( xn ) 2 f ( xn ) f ( xn ) (5.12)
2
Setting, n xn in Eq. 5.12, n 1 (1 2) f ( ) f ( )
n (5.13)
Hence, the Newton – Raphson process has a second order or quadratic convergence.
approximation to the root, as shown in Fig. 5.4. Repeating this process, we approach to the root
quite rapidly. Hence the method consists in replacing the part of the curve between the point
A0 and the x axis by means of the tangent to the curve at A0 .
Geometrically, in Newton-Raphson method, the part of the graph of the function y f ( x)
between the point P ( a, f ( a )) and the x axis is replaced by a tangent to the curve at the
point at each step in the approximation process.
This method is very useful for approximating isolated roots.
The Newton-Raphson method fails if f ( x ) is difficult to compute or vanishes in a
neighbourhood of the desired root. In such cases, the Regula-Falsi method should be used.
The Newton-Raphson method is widely used since in a neighbourhood of the desired root, it
converges more rapidly than the bisection method or the Regula-Falsi method.
If the starting value a is not close enough to the desired root, the method may give a false
root or may not converge.
If f ( x0 ) f ( x0 ) is not sufficiently small, this method does not work. Also if it work, it works
faster.
Example 5.7 [CS-1996, AE-2007, IN-2014 (2 mark)]: Newton-Raphson iteration formula for
3
finding c , c 0 is
2 xn3 3 c 2 xn3 3 c 2 xn3 c 2 xn3 c
(a) xn 1 (b) xn 1 (c) xn 1 2
(d) xn 1 2
3 xn2 3 xn2 3xn 3 xn
Solution (c): Let f ( x ) x 3 c then f ( x ) 0 x 3 c 0 x 3 c . So f ( x) 0 3 x 2 , thus
f ( xn ) x3 c 2 x3 c
from Eq. 5.8 we have the iteration formula as, xn 1 xn xn n 2 n 2 .
f ( xn ) 3 xn 3 xn
Example 5.8 [CS-1997 (1 mark)]: The Newton-Raphson method is used to find the root of the
equation x 2 2 0 . If the iterations are started from –1, the iterations will
(a) converge to –1 (b) converge to 2 (c) converge to 2 (d) not converge
Solution (c): f ( x ) x 2 2 f ( x ) 2 x , so from Eq. 5.8 we have n xn xn 1
xn2 2 x02 2 12 2 3 0 –1 –1.5
xn 1 xn x1 x0 1 (as x0 1 is
2 xn 2 x0 2( 1) 2 1 –1.5 –1.4166
given). The successive iterates are given in the table. Hence the solution 2 –1.4166 –1.4142
of the equation x 2 2 0 with x0 1 converges to 1.4142 2 . 3 –1.4142 –1.4142
range of . Thus with x0 as initial guess we have the root of the given function as x 1.3 . Now if we
start with x1 as initial guess, then a from x1 meets the line b , which meets the x axis at
x 3 5 0.6 which is the roots of the given function. Similarly with x2 as a guess, then a from
x2 meets the line d , whose equation is given as: Ld : y 1 (0.5 1) (1.55 2.05) ( x 2.05)
Ld : x y 1.05 which meets the x axis at x 1.05 ; as x 1.05 in line b gives y 2.25
which is out of range of y so a from x 1.05 does not meets the line b . But x 1.05 in line c
gives y 0.5 which is in the range of y so a from x 1.05 meets the line c , which meets the
x axis at x 1.3 which is a root of the given function. So using x0 , x1 and x2 respectively as
initial guesses, the roots obtained would be 1.3, 0.6, 1.3, respectively.
Example 5.10 [ME-2000 (5 marks)]: Subjective Question: Estimate the root of the equation
(e x x ) 0 to four decimal accuracy by employing the Newton-Raphson method starting with an
initial guess of x0 0 .
Solution: Let f ( x ) e x x f ( x) e x 1 , so from Eq. 5.8 we have n xn xn 1
xn xn xn xn
f ( xn ) e x x e e e ( xn 1) 0 0 0.5
xn 1 xn xn x n n xn
xn 1 x
.
f ( xn ) n
e 1 1 e 1 e n 1 0.5 0.5663
The successive iterates are given in the table. So the solution of 2 0.5663 0.5671
(e x x ) 0 is x 0.5671 . 3 0.5671 0.5671
[Similar questions were also asked in [EC-2008, EE-2008 (2 marks)]
Statement for Linked Answer Questions 5.11 & 5.12: Given a 0 , we wish to calculate its
reciprocal value 1 a by using Newton Raphson method for f ( x ) 0 .
Example 5.11 [CE-2005 (2 marks)]: The Newton-Raphson algorithm for the function will be
1 a a a
(a) xk 1 xk (b) xk 1 xk xk2 (c) xk 1 2 xk axk2 (d) xk 1 xk xk2
2 xk 2 2
Solution: In order to apply Newton-Raphson method to find the reciprocal of a , it is necessary to
find a function f ( x ) which has a zero at x 1 a . The obvious such function is f ( x ) x (1 a ) , but
the Newton–Raphson iteration for this is unhelpful since it cannot be computed without already
knowing the reciprocal of a ; also multiple iterations for refining reciprocal are not possible since
higher order derivatives do not exist for f ( x ) x (1 a ) . A function which does work is
f ( x ) (1 x ) a , for which the Newton–Raphson iteration gives
xk 1 xk f ( xk ) f ( xk ) xk {(1 xk ) a} {1 xk2 } xk (1 axk ) xk 2 xk axk2 .
Example 5.12 [CE-2005 (4 marks)]: For a 7 and (a) 0.11, 0.1299 (b) 0.12, 0.1392
starting with x0 0.2 , the first two iteration will be (c) 0.12, 0.1416 (d) 0.13, 0.1428
Solution (b): Using the result of previous problem with a 7 and x0 0.2 , we have
x1 2 x0 ax02 2
2(0.2) 7(0.2) 0.12 ; x2 2 x2 ax22 2
2(0.12) 7(0.12) 0.1392 .
Example 5.13 [CS-2007 (2 marks)]: Solution (a): The given iterative formula converges
Consider the series xn 1 ( xn 2) 9 (8 xn ) , at n and at that point xn 1 xn x (say); then
x0 0.5 obtained from the Newton-Raphson x 9 18 9 3
we have x x2 x 1.5 .
method. The series converges to 2 8x 8 4 2
(a) 1.5 (b) 2 (c) 1.6 (d) 1.4 So the given series converges to x 1.5 .
Example 5.18 [AE-2009 (2 marks), CE-2011 (1 mark)]: The square root of a number N is to be
obtained by applying the Newton Raphson iterations to the equation x 2 N 0 . If i denotes the
iteration index, the correct iterative scheme will be
1 N 1 2 N 1 N2 1 N
(a) xi 1 xi (b) xi 1 xi 2 (c) xi 1 xi (d) xi 1 xi
2 xi 2 xi 2 xi 2 xi
Solution (a): f ( x) x 2 N f ( x ) 2 x , so from Eq. 5.8 we have the iteration formula as,
f ( xi ) x 2 N xi2 N 1 N
xi 1 xi xi i xi .
f ( xi ) 2 xi 2 xi 2 xi
[Similar questions were also asked in CS-1995, EE-2009, IN-2007, ME-1999 (2 marks)]
Example 5.19 [XE-2009 (1 mark)]: The root of ax b 0 ( a , b constants), can be found by the
Newton-Raphson method with a minimum of
(a) 1 iteration (b) 2 iterations (c) 3 iterations (d) an undeterminable number of iterations
Solution (a): As f ( x ) ax b f ( x) a . So from Eq. 5.9 we have the iteration formula as,
xn 1 xn ( axn b) a b a . Hence we can see that xn 1 b a for any choice of x0 . So the
root of ax b 0 , can be found by the Newton-Raphson method with a minimum of one iteration.
Example 5.20 [AE-2011 (2 marks)]: Consider the function f ( x) x sin x . The Newton-Raphson
iteration formula to find the root of the function starting from an initial guess x (0) at iteration k is
(a) x ( k 1) {sin x ( k ) x ( k ) cos x ( k ) } {1 cos x ( k ) } (b) x ( k 1) {sin x ( k ) x ( k ) cos x ( k ) } {1 cos x ( k ) }
(c) x ( k 1) {sin x ( k ) x ( k ) cos x ( k ) } {1 cos x ( k ) }
(d) x ( k 1) {sin x ( k ) x ( k ) cos x ( k ) } {1 cos x ( k ) }
f ( xk ) x sin xk sin xk xk cos xk
Solution (a): From Eq. 5.8, we have xk 1 xk xk k .
f ( xk ) 1 cos xk 1 cos xk
Example 5.21 [EC-2011 (2 mark)]: A numerical solution of the equation f ( x ) x x 3 0 can
be obtained using Newton-Raphson method. If the starting value is x 2 for the iteration, the value
of x that is to be used in the next step is
(a) 0.306 (b) 0.739 (c) 1.694 (d) 2.306
Solution (c): Let f ( x ) x x 3 f ( x ) 1 {1 (2 x )} , so from Eq. 5.8 we have the iteration
f ( xn ) x xn 3 2 x x 2 xn 6 xn xn 6 xn
formula as, xn 1 xn xn n xn n n . So
f ( xn ) 1 1 2 xn 2 xn 1 2 xn 1
x0 6 x0 2 6 2 6.485
with x0 2 as initial guess, we have x1 1.694 .
2 x0 1 2 2 1 3.828
Example 5.22 [IN-2011 (2 marks)]: The extremum (minimum or maximum) point of a function
f ( x ) is to be determined by solving df ( x) dx 0 using the Newton Raphsons method. Let
f ( x ) x 3 6 x and x0 1 be the initial guess of x . The value of x after two iterations ( x2 ) is
(a) 0.0141 (b) 1.4142 (c) 1.4167 (d) 1.5000
Solution (c): The extremum point of a function is determined by equating its derivative to zero. So
our function g ( x) f ( x ) 3 x 2 6 g ( x ) 6 x , so from Eq. 5.8 we have the iteration formula as,
g ( xn ) 3x2 6 3x02 6 3(12 ) 6
xn 1 xn xn n . Given x0 1 x1 x0 x1 1 1.5 . So,
g ( xn ) 6 xn 6 x0 6(1)
x2 x1 (3x12 6) (6 x1 ) 1.5 [3(1.5 2 ) 6] [6(1.5)] 1.4167 .
Example 5.23 [XE-2011 (2 marks)]: Suppose xn is the nth iterated value while finding the positive
square root of 7 by the Newton-Raphson method with a positive initial guess x0 ( 7 ) . If
en 7 xn for n 1 , then
(a) en 1 en (2 xn2 ) (b) en 1 7 en2 (2 xn ) (c) en 1 en2 7 (d) en 1 en2 (2 xn )
1 7
Solution (d): w.r.t problem 5.17, with N 7 , xn 1 xn . As en 7 xn for n 1 , then
2 xn
1 7 x 7 2 7 xn xn2 ( 7 ) 2 ( 7 xn ) 2 en2
en 1 7 xn 1 7 xn 7 n
2 xn 2 2 xn 2 xn 2 xn 2 xn
Example 5.24 [CH-2012 (2 marks)]: The Newton-Raphson method is used to find the roots of the
equation f ( x ) x cos x , 0 x 1 . If the initial guess for the root is 0.5, then the value of x after
the first iteration is
(a) 1.02 (b) 0.62 (c) 0.55 (d) 0.38
Solution: f ( x ) x cos x f ( x ) 1 sin x , so from Eq. 5.8 we have the iteration formula as,
f ( xn ) x cos xn xn sin xn cos xn
xn 1 xn xn n . So with x0 0.5 as initial guess, we
f ( xn ) 1 sin xn 1 sin xn
have xn 1 { (0.5) sin( 2) cos( 2)} {1 sin( 2)} { (0.5)} (1 ) {2(1 )} 0.38 .
Example 5.25 [EE-2013 (2 marks)]: When the Newton-Raphson method is applied to solve the
equation f ( x ) x 3 2 x 1 0 , the solution at the end of the first iteration with the initial guess value
as x0 1.2 is
(a) –0.82 (b) 0.49 (c) 0.705 (d) 1.69
3 2
Solution: f ( x ) x 2 x 1 f ( x ) 3 x 2 from Eq. 5.8 we have the iteration formula as,
f ( xk ) x3 2 x 1 2 xk3 1
xk 1 xk xk k 2 k 2 . So with x0 1.2 as initial guess we have
f ( xk ) 3 xk 2 3 xk 2
x1 (2 x03 1) (3 x02 2) {2(1.2)3 1} {3(1.2) 2 2} 0.705 .
[Similar questions were also asked in CH-2010 (1 mark), ME-2005, CE-2007, MT-2010, TF-
2008 (2 marks)]
Solution: Let f ( x ) e x 1 f ( x ) e x , so from Eq. 5.8, we have the iterative formula as,
x x x
f ( xn ) e n 1 x e n e n 1
xn 1 xn xn x n x
. So with x0 1 , we have
f ( xn ) en en
x x x x x
x0 e 0 e 0 1 (1)e1 e1 1 1 e 1 1 x1e 1 e 1 1
x1 x
and thus x2 x1 x
x
e0 e1 e e1 e1
1e 1e
(1 e)e e 1 1 1 1
x2 1e
1 1 e 0.368 1 0.06 . Now as x 0 is the root of the
e e e (2.71)0.368
nd
equation f ( x ) 0 . So absolute error at 2 iteration is 0 0.06 0.06 .
Example 5.28 [MA-2014 (2 marks)]: Using the Newton-Raphson method with the initial guess
(0)
x 6 , the approximate value of the real root of x log10 x 4.77 , after the second iteration, is ……
log10 x
Solution: As log10 x (log10 e) (log10 e)(log e x ) 0.434 ln x . So let
log10 e
f ( x ) 0.434 x ln x 4.77 f ( x ) 0.434(ln x 1) 0.434 ln x 0.434 . Thus we have the iterative
f ( xn ) 0.434 xn ln xn 4.77 0.434 xn 4.77
formula as, xn 1 xn xn xn 1 . So with
f ( xn ) 0.434 ln xn 0.434 0.434 ln xn 0.434
0.434 x0 4.77 0.434(6) 4.77
initial guess x0 6 , we have x1 6.086 ; so
0.434 ln x0 0.434 0.434 ln 6 0.434
0.434 x1 4.77 0.434(6.086) 4.77
x2 6.085 .
0.434 ln x1 0.434 0.434 ln 6.086 0.434
Example 5.29 [ME-2014 (2 marks)]: The real root of the equation 5 x 2 cos x 1 0 (up to two
decimal accuracy) is …………….
Solution: Let f ( x ) 5 x 2 cos x 1 f ( x ) 5 2 sin x , so using Newton-Raphson method the
iterative formula which gives the solution of f ( x) 0 is
f ( xn ) 5 x 2 cos xn 1 2 xn sin xn 2 cos xn 1
xn 1 xn xn n .
f ( xn ) 5 2 sin xn 5 2 sin xn
Now as f (0) 2 1 3 0 and f (1) 5 2 cos(1c ) 1 2.92 0 ; so by n xn xn 1
intermediate value theorem, one of the roots of f ( x ) 0 is lying between 0 1 0.56
[0,1] . Now with x0 1 as initial guess, the successive iterates up to two 1 0.56 0.54
decimal accuracy are given below: 2 0.54 0.54
Example 5.30 [PI-2014 (2 marks)]: If the equation sin x x 2 is solved by Newton Raphson’s
method with the initial guess of x 1 , then the value of x after 2 iterations would be ………
Solution: Let f ( x ) sin x x 2 f ( x ) cos x 2 x , so using Newton-Raphson method the iterative
formula which gives the solution of f ( x) 0 is
f ( xn ) sin xn xn2 x cos xn xn2 sin xn
xn 1 xn xn n . So with x0 1 as initial guess, we
f ( xn ) cos xn 2 xn cos xn 2 xn
x0 cos x0 x02 sin x0 cos1c 1 sin1c
have x1 c
0.89 ; similarly
cos x0 2 x0 cos1 2
2
x1 cos x1 x1 sin x1 (0.89) cos 0.89c 0.89 2 sin 0.89c
x2 c
0.87 .
cos x1 2 x1 cos 0.89 2(0.89)
Example 5.31 [CE-2016 (1 mark)]: Newton-Raphson method is to be used to find root of equation
x
3 x e sin x 0 . If the initial trial value for the root is taken as 0.333, the next approximation for
the root would be _____. (note: answer up to three decimal)
Solution: The given function is f ( x) 3 x e x sin x f ( x) 3 e x cos x . So with initial
approximate root x0 0.333 , the next approximated roots is find by
f ( x0 ) 3(0.333) e 0.333 sin(0.333)
x1 x0 0.333 0.333
0.360 .
f ( x0 ) 3e cos(0.333)
[Similar question was also asked in ME-2016 (1 mark)]
Example 5.32 [ME-2016 (1 mark)]: Solve the equation x 10 cos( x) using the Newton-Raphson
method. The initial guess is x 4 . The value of the predicted root after the first iteration, up to
second decimal, is _____.
Solution: The given function is f ( x) x 10cos x f ( x) 1 10sin x . So with initial approximate
root x0 4 , the next approximated roots is given by
f ( x0 ) ( 4) 10 cos( 4) 6.285
x1 x0 1.5644 .
f ( x0 ) 4 1 10 sin( 4) 4 8.071
[Similar question was also asked in AE-2016, MT-2016, PI-2016 (2 marks)]
The Secant Method: Newton’s method is very powerful but has the disadvantage that the
derivative f may sometimes be a far more difficult expression than f itself and its evaluation
therefore computationally expensive. This situation suggests the idea of replacing the derivative with
the difference quotient
f ( xn ) f ( xn 1 )
f ( xn ) (5.14)
xn xn 1
xn xn 1
Hence the Newton – Raphson formula becomes, xn 1 xn f ( xn ) (5.15)
f ( xn ) f ( xn 1 )
Geometrically, we intersect the x axis at xn 1 with the secant of f ( x )
passing through Pn 1 and Pn in Fig. 5.5. We need two starting values
x0 and x1 . Evaluation of derivatives is now avoided. It can be shown
that convergence is more rapid than linear, almost quadratic like
Newton – Raphson method. The algorithm is similar to that of
Newton’s method.
It should be noted that this formula requires two initial
approximations to the root. Figure 5.5: Secant Method
x f ( xn ) xn f ( xn 1 )
It is not good to write Eq. 5.15 as, xn 1 n 1 because this may lead to loss of
f ( xn ) f ( xn 1 )
significant digits if xn and xn 1 are about equal.
Example 5.33 [MT-2013 (2 marks)]: Applying the secant method, the first approximation to the root
of f ( x ) 1 ln x ( x 2) , starting with function values at x 0.3 and x 0.4 is ……
Solution: The two initial approximations are x0 0.3 and x1 0.4 . So from Eq.5.15,
xn xn 1
xn 1 xn f ( xn ) . Thus with n 1, we have the first approximation for the root of
f ( xn ) f ( xn 1 )
x1 x0
f ( x ) as x2 x1 f ( x1 ) . As f ( x0 ) f (0.3) 1 ln 0.3 (0.3 2) 0.054 and
f ( x1 ) f ( x0 )
0.4 0.3
f ( x1 ) f (0.4) 1 ln 0.4 (0.4 2) 0.284 . So x2 0.4 0.284 0.316 .
0.284 0.054
The Rate of Convergence for the iterative method discussed in previous sections is given as:
The Bisection Method: Order of Convergence k 1 , linear convergence
The method of False Position: Order of Convergence k 1 , linear convergence.
Fixed – Point method: Order of Convergence k 1 , linear convergence.
Newton – Raphson method: Order of Convergence k 2 , quadratic convergence [This point
was asked in AE-2010 (2 marks)].
The Secant method: Order of Convergence k 1.618 .
5.1.5 Interpolation
Consider the following problem: Given the values of a known function y f ( x) at a sequence of
ordered points x0 , x1 , , xn , find f ( x ) for arbitrary x . When x0 x xn , the problem is called
interpolation. When x x0 or x xn , the problem is called extrapolation. With yi f ( xi ) , the
problem of interpolation is basically one of drawing a smooth curve through the known points
( x0 , y0 ), ( x1 , y1 ), , ( xn , yn ) . So what kind of function f ( x ) should one choose? A polynomial is a
common choice for an interpolating function because polynomials are easy to evaluate, differentiate,
and integrate relative to other choices such as a trigonometric and exponential series. Polynomial
interpolation involves finding a polynomial of order n that passes through the n 1 points. The
methods of interpolation that we discuss are Lagrangian interpolation method and the Newton’s
divided difference polynomial method.
Example 5.34 [MN-2012 (1 mark)]: Assuming sin(1) 0.841 and sin(3) 0.141 , the Lagrangian
linear interpolation polynomial, for the function f ( x ) sin x defined on the interval [1, 3] and
passing through the end points of the intervals, is
(a) 0.35 x 1.19 (b) 3.05 x 11.92 (c) 35.00 x 119.10 (d) 40.50 x 219.19
Solution (a): From Eq. 5.17, we have f ( x ) sin x , x0 1 , x1 3 , f ( x0 ) 0.841 and
f ( x1 ) 0.141 . So p1 ( x ) ( x 3) (1 3) 0.841 ( x 1) (3 1) 0.141 1.191 .35 x .
General Lagrange Interpolation Polynomial: Generalising Eq. 5.16 and 5.17, for n , i.e.
n n
f ( x) pn ( x) k 0 Lk ( x) f k k 0 {lk ( x)} {lk ( xk )} f k (5.18)
where, Lk ( xk ) 1 and Lk 0 at the other nodes; and the Lk are independent of the function f to be
interpolated; l0 ( x) ( x x1 )( x x2 ) ( x xn ) ;
lk ( x ) ( x x0 ) ( x xk 1 )( x xk 1 ) l0 ( x) ( x x1 )( x x2 ) ( x xn ), 0k n;
ln ( x ) ( x x0 )( x x1 ) ( x xn 1 ) ; pn ( xk ) fk ; and lk ( x j ) 0 if j k .
Error Estimate: If f is itself a polynomial of degree n (or less), it must coincide with pn because
the n 1 data ( x0 , f 0 ), , ( xn , f n ) determine a polynomial uniquely, so the error is zero. Now the
special f has its ( n 1)th derivative identically zero. This makes it plausible that for a general f its
( n 1)th derivative f ( n 1) should measure the error n ( x ) f ( x ) pn ( x ) . It can be shown that this is
true if f ( n 1) exists and is continuous. Then, with a suitable t between x0 and xn , we have,
n ( x ) f ( x ) pn ( x) ( x x0 )( x x1 ) ( x xn ) 1 ( n 1)! f ( n 1) (t ) (5.19)
Thus n ( x) is 0 at the nodes and small near them, because of continuity. The product
( x x0 ) ( x xn ) is large for x away from the nodes. This makes interpolation at an x will be best
if we choose nodes on both sides of that x . Also we get error bounds by taking the smallest and
largest value of f ( n 1) (t ) in Eq. 5.19 on the interval x0 t xn .
Newton’s Divide Difference Interpolation: Suppose that Pn ( x ) is the nth Lagrange polynomial
that agrees with the function f at the distinct numbers x0 , x1 , , xn . Although this polynomial is
unique, there are alternate algebraic representations that are useful in certain situations. The divide
differences of f with respect to x0 , x1 , , xn are used to express Pn ( x ) in the form,
Pn ( x ) a0 a1 ( x x0 ) a2 ( x x0 )( x x1 ) an ( x x0 )( x x1 ) ( x xn 1 ) (5.20)
for appropriate constants a0 , a1 , , an . To determine the first of these constants, a0 , we have to
evaluate Pn ( x ) (from Eq. 5.20) at x x0 , which gives a0 Pn ( x0 ) f ( x0 ) . Similarly, when P ( x ) is
evaluated at x1 , the only nonzero terms in the evaluation of Pn ( x1 ) are the constant and linear terms,
f ( x0 ) a1 ( x1 x0 ) Pn ( x1 ) f ( x1 ) a1 { f ( x1 ) f ( x0 )} {x1 x0 } (5.21)
We now introduce the divided – difference notation. The zeroth divided difference of the function f
with respect to xi , denoted as f [ xi ] , which is simply the value of f at xi , i.e.,
f [ xi ] f ( xi ) (5.22)
The remaining divided difference are defined recursively; the first divided difference of f with
respect to xi and xi 1 is denoted as f [ xi , xi 1 ] and defined as,
f [ xi , xi 1 ] { f [ xi 1 ] f [ xi ]} {xi 1 xi } (5.23)
The second divided difference, f [ xi , xi 1 , xi 2 ] , is defined as,
f [ xi , xi 1 , xi 2 ] { f [ xi 1 , xi 2 ] f [ xi , xi 1 ]} { xi 2 xi } (5.24)
th
Similarly, after the ( k 1) divided differences, f [ xi , xi 1 , xi 2 , , xi k 1 ] and f [ xi 1 , xi 2 , , xi k ]
have been determined. So the k th divided difference relative to xi , xi 1 , xi 2 , , xi k is given as,
f [ xi , xi 1 , xi 2 , , xi k ] { f [ xi 1 , xi 2 , , xi k ] f [ xi , xi 1 , xi 2 , , xi k 1 ]} {xi k xi } (5.25)
th
The process ends with the single n divided difference,
f [ x0 , x1 , x2 , , xn ] { f [ x1 , x2 , , xn ] f [ x0 , x1 , , xn 1 ]} {xn x0 } . Because of Eq. 5.21, we can
write a1 [ x0 , x1 ] , just as a0 can be expressed as a0 f ( x0 ) f [ x0 ] . Hence the interpolating
polynomial in Eq. 5.20 is given as,
Pn ( x ) f [ x0 ] f [ x0 , x1 ]( x x0 ) a2 ( x x0 )( x x1 ) an ( x x0 )( x x1 ) ( x xn 1 ) (5.26)
As might be expected from the evaluation of a0 and a1 , the required constants are
ak f [ x0 , x1 , x2 , xk ] for each k 0,1, , n . So Pn ( x ) can be re-written in a form called Newton’s
Divided Difference, as
n
Pn ( x ) f [ x0 ] k 1 f [ x0 , x1 , , xk ]( x x0 ) ( x xk 1 ) (5.27)
The value of f [ x0 , x1 , , xk ] is independent of the order of the numbers x0 , x1 , , xk .
Example 5.36: For x, f ( x ) (1, 0.765), (1.3, 0.620), (1.6, 0.455), (1.9, 0.281), (2.2, 0.110) , compute
the divided difference table for the given data and construct the interpolating polynomial.
Solution: The first divided difference involving x0 1 and x1 1.3 is given as,
f [ x0 , x1 ] { f [ x1 ] f [ x0 ]} { x1 x0 } {0.620 0.765} {1.3 1.0} 0.483 .
The remaining first divided differences are found in a similar manner and are shown in the fourth
column of the Table. The second divided difference involving x0 , x1 and x2 is
f [ x0 , x1 , x2 ] { f [ x1 , x2 ] f [ x0 , x1 ]} {x2 x0 } {0.548 ( 0.483)} {1.6 1.0} 0.108 . The
remaining second divided differences are found in a similar manner and are shown in the fifth column
of the Table. The third divided difference involving x0 , x1 , x2 and x3 is
f [ x0 , x1 , x2 , x3 ] { f [ x1 , x2 , x3 ] f [ x0 , x1 , x2 ]} { x3 x0 } {0.049 ( 0.108)} {1.9 1.0} 0.065 ;
the other is evaluated similarly. The fourth divided difference involving x0 , x1 , x2 , x3 and x4 is
f [ x0 , x1 , x2 , x3 , x4 ] { f [ x1 , x2 , x3 , x4 ] f [ x0 , x1 , x2 , x3 ]} {x4 x0 } {0.068 0.065} {2.2 1} 0.002
i xi f [ xi ] f [ xi 1 , xi ] f [ xi 2 , xi 1 , xi ] f [ xi 3 , , xi ] f [ xi 4 , , xi ]
0 1.0 0.765
– 0.483
1 1.3 0.620 – 0.108
– 0.548 0.065
2 1.6 0.455 – 0.049 0.002
– 0.578 0.068
3 1.9 0.281 0.011
– 0.571
4 2.2 0.110
The coefficients of the Newton forward divided – difference form of the interpolating polynomial are
along the diagonal in the table. This polynomial is,
P4 ( x) 0.765 0.483( x 1.0) 0.108( x 1.0)( x 1.3) 0.065( x 1.0)( x 1.3)( x 1.6)
0.002( x 1.0)( x 1.3)( x 1.6)( x 1.9)
Simplified form of Newton’s divided difference formula: Newton divided difference formula
can be expressed in a simplified form when the nodes are arranged consecutively with equal spacing.
In this case, we introduce the notation h xi 1 xi , for each i 0,1, , n 1 and let x x0 sh .
Then the difference x xi is x xi ( s i) h . So Eq. 5.27 can be written as
Pn ( x ) Pn ( x0 sh) f [ x0 ] shf [ x0 , x1 ] s ( s 1) h 2 f [ x0 , x1 , x2 ]
s ( s 1) ( s n 1)h n f [ x0 , x1 , , xn ]
n
Pn ( x ) f [ x0 ] k 1 s ( s 1) ( s k 1)h k f [ x0 , x1 , , xk ] . Using binomial-coefficient notation,
n s ( s 1) ( s k 1)
Pn ( x ) Pn ( x0 sh ) f [ x0 ] k 1 s Pk k ! h k f [ x0 , xi , , xk ] , where s Pk (5.28)
k!
Equal Spacing: Newton’s Forward Difference Formula: Newton formula (Eq. 5.27) is valid
for arbitrarily spaced nodes as they may occur in observations. In many applications the x j ’s are
regularly spaced, like, in measurements taken at regular intervals of time; then denoting the distance
by h , we can write the x j ’s as x0 , x1 x0 h, x2 x0 2h, , xn x0 nh . The first forward
difference of f at x j is given as, f j f j 1 f j ; the second forward difference of f at x j is given
2
as, f j f j 1 f j ; and continuing in this way, the k th forward difference of f at x j is given as
k f j k 1 f j 1 k 1 f j , k 1, 2, . Now forward – difference formula is constructed by making
use of the forward difference notation ‘ ’.
In Eq. 5.27 we finally set x x0 rh x x0 rh ; similarly x x1 ( r 1)h , and so on. With this
and Eq. 5.29, Eq. 5.27 becomes Newton’s forward difference interpolation formula.
n r (r 1) 2 r ( r 1) ( r n 1) n
f ( x ) pn ( x ) s 0 r Ps s f 0 f 0 r f 0 f0 f0 (5.30)
2! n!
r (r 1)(r 2) ( r s 1)
where r ( x x0 ) h , the binomial coefficients, r Ps ( s 0, integer) .
s!
Example 5.37 [ME-2004 (2 marks)]: The values of a function f ( x ) are tabulated below:
x 0 1 2 3 Using Newton’s forward difference formula, the cubic
f ( x) 1 2 1 10 polynomial that can be fitted to the above data is
(a) 2 x 3 7 x 2 6 x 2 (b) 2 x3 7 x 2 6 x 2 (c) x 3 7 x 2 6 x 1 (d) 2 x 3 7 x 2 6 x 1
Solution (d): From the given i xi f [ xi ] f [ xi 1 , xi ] f [ xi 2 , xi 1 , xi ] f [ xi 3 , xi 2 , xi 1 , xi ]
data, we note that the x j ’s
0 0 1
are regularly spaced, with 1
the spacing distance h 1,
1 1 2 –1
then we have x0 0 ,
–1 2
x1 x0 h 1 , ,
2 2 1 5
x3 x0 3h 3 . So the
9
difference table is given as:
3 3 10
So from Eq. 5.29, we have f ( x0 ) 1! h f [ x0 , x1 ] (1!) 1 (1) 1 ;
f ( x0 ) 2! h f [ x0 , x1 , x2 ] (2!) 1 ( 1) 2 ; f ( x0 ) 3! h 3 f [ x0 , x1 , x2 , x3 ] (3!) 13 2 12 .
2 2 2 3
r ( r 1) r ( r 1)( r 2)
Thus from Eq. 5.39, f ( x ) pn ( x) f ( x0 ) r f ( x0 ) 2 f ( x0 ) 3 f ( x0 ) ,
2! 3!
x ( x 1) x ( x 1)( x 2)
where x x0 rh x 0 r (1) r x , so f ( x ) 1 x(1) ( 2) (12)
2! 3!
f ( x) 1 x x 2 x 2 x3 6 x2 6 x 2 x3 7 x 2 6 x 1 .
Equal Spacing: Newton’s Backward Difference Formula: If the interpolating nodes are
recorded from last to first as xn , xn 1 , , x0 , we can write the interpoatory formula as,
Pn ( x ) f [ xn ] f [ xn , xn 1 ]( x xn ) f [ xn , xn 1 , xn 2 ]( x xn )( x xn 1 )
(5.31)
f [ xn , , x0 ]( x xn )( x xn 1 ) ( x x1 )
We define the first backward difference of f at x j by f j f j f j 1 ; the second backward
2
difference of f at x j is given as, f j f j f j 1 ; and continuing in this way the k th backward
difference of f at x j is given as,
k f j k 1 f j k 1 f j 1 , ( k 1, 2, ) (5.32)
A formula similar to Eq. 5.30 but involving backward difference is Newton’s backward difference
interpolation formula is given as,
n r ( r 1) 2 r (r 1) ( r n 1) n
f ( x ) p n ( x ) s 0 ( r s 1) Ps s f 0 f 0 r f 0 f0 f 0 (5.33)
2! n!
where, x x0 rh , r ( x x0 ) h .
Newton’s Central Differences: The Newton forward- and backward-difference formulas are not
appropriate for approximating f ( x ) when x lies near the centre of the table because neither will
permit the highest-order difference to have x0 close to x . A number of divided-difference formulas
are available for this case, each of which has situations when it can be used to maximum advantage.
These methods are known as centred-difference formulas, x1 f 1
which are third notation for differences. The first central
f 1 2
difference of f ( x ) at x j is defined by f j f j 1 2 f j 1 2 ;
x0 f0 2 f0
and the k th central difference of f ( x ) at x j is given as,
k k 1 k 1 f1 2 3 f1 2
f j f j 1 2 f j 1 2 , j 2, 3, (5.34)
Thus in this notation a difference table, for example, for
x1 f1 2 f1
f 1 , f0 , f1 , f 2 looks as given in the table. Central f3 2
differences are used in numeric differentiation, differential x2 f2
equation and central interpolation formulas.
We can use Taylor polynomials to derive the accuracy of the Forward, Backward and Central
difference formulas. For e.g., the usual form of the Taylor polynomial with remainder is
h2
f ( x h ) f ( x) h f ( x ) f (c) , where c is some unknown number between x and x h .
2
f ( xi 1 ) f ( xi ) h
Letting x xi , x h xi 1 and solving for f ( xi ) leads to f ( xi ) f (c ) . Notice
h 2
that the quotient in this equation is exactly the forward difference formula. Thus the error of the
h
forward difference is f (c) which means it is O( h) . Replacing h in the above calculation by h
2
gives the error for backward difference formula, i.e., it is also O( h) . For the central difference, the
error can be found from the third degree Taylor polynomial with remainder
h2 h3
f ( xi 1 ) f ( xi h) f ( xi ) h f ( xi ) f ( xi ) f (c1 ) and
2 3!
h2 h3
f ( xi 1 ) f ( xi h) f ( xi ) h f ( xi )
f ( xi ) f (c2 ) , where xi c1 xi 1 and xi 1 c2 xi
2 3!
. Subtracting these two equations and solving for f ( xi ) leads to
f ( xi 1 ) f ( xi 1 ) h 2 f (c1 ) f (c2 )
f ( xi ) . This shows that the error for the central difference
2h 3! 2
formula is O ( h 2 ) . Thus, central differences are significantly better and so. It is best to use central
differences whenever possible.
Example 5.38 [CS-1996 (1 mark)]: The formula used to compute an approximation for the second
derivative of a function f at a point x0 is
f ( x0 h) f ( x0 h ) f ( x0 h) f ( x0 h )
(a) (b)
2 2h
f ( x0 h) 2 f ( x0 ) f ( x0 h ) f ( x0 h) 2 f ( x0 ) f ( x0 h )
(c) (d)
h2 h2
h2
Solution (d): From 2nd degree Taylor polynomial, we have f ( x0 h) f ( x0 ) h f ( x0 ) f ( x0 )
2
h2
…(i); and f ( x0 h) f ( x0 ) h f ( x0 ) f ( x0 ) …(ii). Adding (i) and (ii), we get
2
f ( x0 h) 2 f ( x0 ) f ( x0 h)
f ( x0 h) f ( x0 h ) 2 f ( x0 ) h 2 f ( x0 ) f ( x0 ) .
h2
df ( x ) f ( x0 h) f ( x0 h)
estimated with h 0.03 using the central difference formula is
dx x x0 2h
3
2 10 . The values of x0 and f ( x0 ) are 19.78 and 500.01, respectively. The corresponding error in
the central difference estimate for h 0.02 is approximately
(a) 1.3 10 4 (b) 3.0 104 (c) 4.5 104 (d) 9.0 104
h 2 f (c1 ) f (c2 )
Solution (d): As the error in central difference formula is E , where
3! 2
c1 ( x0 , x0 h) and c2 ( x0 h, x0 ) . So for h 0.03 , we have
3 (0.03) 2 f (c1 ) f (c2 ) (0.02) 2 f (c1 ) f (c2 )
E1 2 10 . For h 0.02 , we have E2 . So
3! 2 3! 2
E2 (0.02) 2 4
2
E 2 2 10 3 8.89 10 4 9 10 4 .
E1 (0.03) 9
Example 5.40 [ME-2014 (1 mark)]: The best approximation of the minimum value attained by
e x sin(100 x ) for x 0 is …………….
Solution: y e x sin(100 x ) y e x {100 cos(100 x ) sin(100 x)} , and
y e x {9999 sin(100 x ) 200 cos(100 x)} . For maxima and minima, y 0 , since e x 0 , so
{100 cos(100 x ) sin(100 x )} 0 tan(100 x ) 100 100 x tan 1 (100) 89.43o or 269.43o
x 0.8943o or 2.6943o ; and y x 0.8943o 0 and y x 2.6943o 0 . So x 0.8943o is the point of
o
maxima and x 2.6943 is the point of minima and thus the minimum value will be
y e 2.6943( 180) sin(100 2.6943) 0.954 0.95 . Note that while calculating the power of
exponential we have to convert x 2.6943o in radians. So answer is 0.95 .
Example 5.41 [MT-2008 (1 mark)]: The value of dy dx for the following data set at x 3.5 ,
computed by central difference method is
x 1 2 3 4 5 (a) 3.5 (b) 7
y 0 3 8 15 24 (c) 10.5 (d) 14
dy f (4) f (3) 15 8
Solution: From central difference method, we have 7.
dx x 3.5 43 1
[Similar question was also asked in ME-1999 (1 mark)]
Example 5.43 [MT-2016 (2 marks)]: A hot body cools according to the following equation
dT
cT , where T is the instantaneous temperature at time t , and the constant c 0.05s 1 .
dt
Reduce the differential equation into its finite difference from using forward difference. For
maintaining numerical stability, the maximum value of the time step t (in seconds) is _____.
dT Ti 1 Ti
Solution: Using forward difference f (ti ) Ti . So after converting the given DE into
dt ti 1 ti
Ti 1 Ti
its finite difference from using forward difference, we get cTi , where
ti 1 ti
ti 1 ti t ti 1 ti t value of the time step. So Ti 1 Ti ctTi Ti 1 (1 ct )Ti ; hence
the solution of DE is stable if (1 ct ) 1 1 1 ct 1 0 t (2 c ) , hence the maximum
value of the time step t (in seconds) is 2 c 2 0.05 40 .
Exercise: 5.1
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. The number 0.0008857 when rounded off to three significant digits yields
(a) 0.001 (b) 0.000886 (c) 0.000885 (d) None of these
2. When a number is approximated to n decimal places by chopping off the extra digits, then the
absolute value of the relative error does not exceed
(a) 10 n (b) 10 n 1 (c) 0.5 10 n 1 (d) None of these
3. If e1 and e2 are absolute errors in two numbers n1 and n2 respectively due to rounding or
e1 e2
truncation, then is
n1 n2
(a) e1 e2 (b) e1 e2 (c) e1 e2 (d) e1 e2
3
4. The root of the equation x 6 x 1 0 lies in the interval
(a) (2,3) (b) (3, 4) (c) (3,5) (d) (4, 6)
5. For the equation f ( x ) 0 , if f ( a ) 0, f (b) 0, f (c ) 0 and b c then we will discard the
value of the function f ( x ) at the point
(a) a (b) b (c) c (d) Either a or b or c
x
6. The positive root of the equation e x 3 0 lies in the interval
(a) (0,1) (b) (1, 2) (c) (2,3) (d) (2, 4)
7. If f ( a ) f (b) 0 , then an approximate value of a real root of f ( x ) 0 lying between a and b is
given by
af (b) bf ( a ) bf ( a ) af (b ) af (b) bf ( a)
(a) (b) (c) (d) None of these
ba ba f (b) f ( a )
8. A root of the equation x 3 x 1 0 lies between 1 and 2. Its approximate value as obtained by
applying bisection method 3 times is
(a) 1.375 (b) 1.625 (c) 1.125 (d) 1.25
x
9. The nearest real root of the equation xe 2 0 correct to two decimal places, is
(a) 1.08 (b) 0.92 (c) 0.85 (d) 0.80
3
10. By the false position method, the root of the equation x 9 x 1 0 lies in interval (2, 4) after
first iteration. It is
(a) 3 (b) 2.5 (c) 3.57 (d) 2.47
3
11. The equation x 3 x 4 0 has only one real root. What is its first approximate value as
obtained by the method of false position in (–3, –2)?
(a) –2.125 (b) 2.125 (c) –2.812 (d) 2.812
12. Newton-Raphson method is applicable only when
(a) f ( x ) 0 in the neighborhood of actual root x
(b) f ( x) 0 in the neighborhood of actual root x
(c) f ( x ) 0 in the neighborhood of actual root x
(d) None of these
13. Newton-Raphson processes has a
(a) Linear convergence (b) Quadratic convergence
(c) Cubic convergence (d) None of these
3
14. The real root of the equation x x 5 0 lying between –1 and 2 after first iteration by
Newton-Raphson method is
(a) 1.909 (b) 1.904 (c) 1.921 (d) 1.940
15. The Newton-Raphson method converges fast if f ( ) is ( is the exact value of the root)
(a) Small (b) Large (c) 0 (d) None of these
16. If one root of the equation f ( x ) 0 is near to x0 then the first approximation of this root as
calculated by Newton-Raphson method is the abscissa of the point where the following straight
line intersects the x axis
5.2 Numerical Solution of Definite Integrals, First Order ODEs, and System of
Linear Equations
5.2.1 Numerical Solution of Definite Integrals
b
Numerical integration means the numeric evaluation of integrals I f ( x ) dx where a and b are
a
given and f is a function given analytically by a formula or empirically by a table of values.
Geometrically I is the area under the curve of I between a and b . We know that if f is such that
we can find a differentiable function F whose derivative is f , then we can evaluate I by applying
b
the familiar formula I f ( x ) dx F (b ) F (a ) , where F ( x ) f ( x) .
a
Many times the anti-derivative of the integrand is not easy to obtain. The Newton-Cotes formulas are
an extremely useful and straightforward family of numerical integration techniques. To integrate a
function f ( x ) over some interval [ a, b] , divide it into n equal parts such that f n f ( xn ) and
h (b a ) n . Then find polynomials which approximate the tabulated function, and integrate them
to approximate the area under the curve. To find the fitting polynomials, use Lagrange interpolating
polynomials. The resulting formulas are called Newton-Cotes formulas, or quadrature formulas.
Newton-Cotes formulas may be ‘closed’ if the interval [ x1 , xn ] is included in the fit, ‘open’ if the
points [ x2 , xn 1 ] are used, or a variation of these two. We will consider only closed Newton-Cotes
b
formulas. The basic method involved in approximating a f ( x) dx is called numerical quadrature. It
n b
uses a sum i 0 ai f ( xi ) to approximate a f ( x) dx . The method of quadrature in this section are
based on the interpolation polynomials, discussed in Section 5.1.5. The basic idea is to select a set of
distinct nodes { x0 , x1 , , xn } from the interval [ a, b] . Then integrate the Lagrange interpolating
n
polynomial Pn ( x ) i 0 f ( xi ) Li ( x ) and its truncation error term over [ a, b] to obtain,
( n 1)
b b n b n f ( x)
a f ( x) dx
a
i0 f ( xi ) Li ( x)dx a i0 ( x xi ) ( n 1)!
dx
b n 1 b n ( n 1)
f ( x ) dx i 0 ai f ( xi )
a i0 ( x xi ) f ( x) dx , where ( x ) [a, b] x and
( n 1)! a
b b n
ai Li ( x ) dx , i 0,1, , n . The quadrature formula is therefore a f ( x )dx i 0 ai f ( xi ) with
a
1 b n ( n 1)
error E ( f ) i0 ( x xi ) f ( x) dx . Now, let us consider formulas produced by
( n 1)! a
using first and second Lagrange polynomial with equally – spaced nodes. This gives the Trapezoidal
rule and Simpson’s rule.
Rectangular Rule and Trapezoidal Rule: Numeric integration are obtained by approximating
the integrand f by functions that can be easily integrated. The simplest formula, the rectangular rule,
is obtained if we subdivide the interval of integration a x b into n subintervals of equal length
h (b a ) n and in each subinterval approximate f by the constant f ( x*j ) , the value of f at the
*
midpoint x j of the j th subinterval, as shown in Fig. 5.7 (a). Then f is approximated by a step
function (piecewise constant function), the n rectangles in Fig. 5.7 (a) have the areas
f ( x1* ) h, , f ( x*n ) h and the rectangular rule is given as,
b
I f ( x) dx h[ f ( x1* ) f ( x2* ) f ( xn* )] , where h (b a ) n (5.35)
a
Error Bounds and Estimate for the Trapezoidal Rule: An error estimate for the trapezoidal
rule can be derived from Eq. 5.19 with n 1 by integration as follows. For a single subinterval we
have, f ( x ) p1 ( x ) ( x x0 )( x x1 )(1 2) f (t ) with a suitable t depending on x , between x0 and x1
. Integration over x from a x0 to x1 x0 h gives,
x0 h x0 h
x f ( x) dx (h 2)[ f ( x0 ) f ( x1 )] ( x x0 )( x x0 h)(1 2) f t ( x ) dx . Setting ( x x0 ) v
0 x0
and applying the mean value theorem, which we can use because ( x x0 )( x x0 h) does not change
the sign, we find the right side equals to
h f (t ) h3 h3 f (t ) h3
0 v ( v h ) dv
3 2 2 f (t ) (5.37)
2 12
where t is a (suitable, unknown) value between x0 and x1 . This is the error for the trapezoidal rule
with n 1, often called the local error. Hence the error T of Eq. 5.36 with any n is the sum of such
contributions from the n subintervals; as h (b a ) n , nh3 (b a)3 n 2 , (b a ) 2 n 2 h 2 , thus,
nh3 (b a )3 (b a ) 2
T f (tˆ ) 2
f (tˆ) (5.38) h f (tˆ)
12 12n 12
with (suitable, unknown) tˆ between a and b . Because of Eq. 5.38, the trapezoidal rule (Eq. 5.36) is
also written as,
b 1 1 (b a) 2 ˆ
I f ( x) dx h f ( a) f ( x1 ) f ( x2 ) f ( xn 1 ) f (b) h f (t ) (5.39)
a
2 2 12
Error bounds are now obtained by taking the largest value for f , say, M 1 , and the smallest value,
M 2 , in the interval of integration. Then Eq. 5.38 gives (note that k is negative)
(b a )3 ba 2 n 3
k M 1 T k M 2 , where k 2
(5.40)h h
12n 12 12
We also note that the error term for the Trapezoidal rule involves f , so the Trapezoidal rule
gives the exact result when applied to any function whose second derivative is identically zero,
i.e., the Trapezoidal rule gives the exact result for polynomials of degree up to or equal to one
[This point was asked in CS-2002 (1 mark)].
b
Example 5.44 [CS-1997 (2 marks)]: The trapezoidal method to numerically obtain a f ( x) dx has an
error E bounded by (b a ) 12 h 2 max f ( x ) , x [a, b] , where h is the width of the trapezoids.
Minimum number of trapezoids guaranteed to ensure E 104 in computing ln 7 using f 1 x is
(a) 60 (b) 100 (c) 600 (d) 10000
7
Solution (c): As ln 7 (1 x) dx , so we have a 1 , b 7 and h (b a ) n , where n is the
1
number of intervals. From the given data we have the maximum error of 104 . Thus from Eq. 5.38 we
(b a ) 2 1 1 2
have 104 h max{ f (tˆ)} . Since f ( x ) f ( x) 2 f ( x ) 3 , so in [1, 7] ,
12 x x x
max{ f ( x )} is at x 1 as f ( x ) 2 x 3 which is a decreasing function in [1, 7] . Thus
2 (7 1) 7 1
max{ f (tˆ)} 3
2 . So 104 h 2 2 h 102 10 2 n 600 .
x x 1 12 n
Example 5.45 [CE-2006 (2 marks)]: A 2nd degree polynomial, f ( x ) , has values of 1, 4 and 15 at
2
x 0, 1 and 2, respectively. The integral 0 f ( x )dx is to be estimated by applying the trapezoidal
rule to this data. What is the error (defined as ‘true value – approximate value’) in the estimate?
(a) 4 3 (b) 2 3 (c) 0 (d) 2 3
Solution: From the given data, we have a 0 , b 2 , h 1 and n 2 . So using Trapezoidal rule,
2
IT f ( x ) dx h (1 2) f ( a ) f ( x1 ) (1 2) f (b) 1(1 2) 4 (1 2)(15) 12 , where x1 a h .
0
Now for the exact integral I E , we have to find the function f ( x ) . As f ( x ) is a 2nd degree
polynomial, so f ( x ) ax 2 bx c and it is given that f (0) 1 , f (1) 4 and f (2) 15 ; after
substituting the value and solving for the three unknown from the three equations we get a 4 ,
b 1 and c 1. Thus we have f ( x) 4 x 2 x 1 . So
2
2 2 4 3 x2 32 4 32 32 4
I E (4 x x 1) dx x x 20 . Thus I E IT 12 .
0 3 3
3 2 0 3 2 3
Example 5.46 [ME-2007 (2 marks)]: A calculator has accuracy up to 8 digits after decimal place.
2
The value of 0 sin x dx when evaluated using this calculator by trapezoidal method with 8 equal
intervals, to 5 significant digits is
(a) 0.00000 (b) 1.0000 (c) 0.00500 (d) 0.00025
Solution: From the given data, we have a 0 , b 2 and we have 8 intervals, i.e., n 8 , and so
we have h (b a ) n (2 0) 8 4 . So from Eq. 5.36 with f ( x ) sin x , we have
2
I f ( x ) dx h (1 2) f ( a ) f ( x1 ) f ( x2 ) f ( x7 ) (1 2) f (b ) , where xi a ih . So
0
1 3 5 3 7 1
I
4 2 f (0) f 4 f 2 f 4 f f 4 f 2 f 4 2 f (2 )
1 1 1 1 1 1
I (0) 1 0 1 (0) 0 . So value of I to
4 2 2 2 2 2 2
5 significant digits is I 0.00000 .
Example 5.47 [CS-2008 (2 marks)]: The minimum number of equal length subintervals needed to
2
approximate x
1 xe dx to an accuracy of at least (1 3) 106 using the trapezoidal rule is
(a) 1000e (b) 1000 (c) 100e (d) 100
Solution (a): From the given data we have a 1 , b 2 and h (b a ) n , where n is the number of
intervals. As we have the least accuracy, i.e., maximum error of (1 3) 106 . Thus from Eq. 5.38 we
have (1 3) 106 (b a ) 12 h 2 max{ f (tˆ)} . Since f ( x ) xe x f ( x ) e x (1 x)
f ( x ) e x (2 x ) , so in [1, 2] , max{ f ( x )} is at x 2 as f ( x ) e x (2 x ) which is an
increasing function in [1, 2] . Thus max{ f (tˆ)} e x (2 x) 4e2 . So
x 2
0, t 0 t
u (t ) . So x(t ) e 2t dt . Now using Trapezoidal method with single interval,
0
1, t 0
1 1 1 1 1 e 0.02
x (t ) h f (0) f (t ) x (0.01) h f (0) f (0.01) 0.01 0.0099 .
2 2 2 2 2 2
dx
Example 5.49 [AE-2009 (2 marks)]: The value of the integral 0 1 x sin x evaluated using the
where a 0 , x1 a h 0 ( 2) 2 , b .
1 1 1 1 1 1
So I f (0) f f ( ) I (1) 1.414 .
2 2 2 2 2 2 2 ( 2) 2 1
[Similar questions were also asked in AG-2012, AE-2007, CH-2009, TF-2008 (2 marks)]
1 x2
Example 5.50 [PI-2011 (2 marks)]: The value of 0 e dx , using trapezoidal rule for 10 trapezoids,
is equal to
(a) 0.6778 (b) 0.7165 (c) 0.6985 (d) 0.7462
Solution: From the given data, we have a 0 , b 1 and n 10 , so we have
2
h (b a ) n (1 0) 10 0.1 . So from Eq. 5.36 with f ( x) e x , we have
1
I f ( x) dx h (1 2) f (a ) f ( x1 ) f ( x2 ) f ( x9 ) (1 2) f (b) , where xi a ih i(0.1) . So
0
0.12
I (1 10) (1 2) f (0) f (0.1) f (0.2) f (0.9) (1 2) f (1) . As f (0) 1 , f (0.1) e 0.99
2 2 2 2
, f (0.2) e 0.2 0.96 , f (0.3) e 0.3 0.91 , f (0.4) e 0.4 0.85 , f (0.5) e 0.5 0.78 ,
2 2 2 2
f (0.6) e 0.6 0.70 , f (0.7) e 0.7 0.61 , f (0.8) e 0.8 0.53 , f (0.9) e 0.9 0.44 ,
2
f (1) e 1 0.37 . Thus we have
I (1 10) (1 2)(1) 0.99 0.96 0.91 0.85 0.78 0.70 0.61 0.53 0.44 (1 2)(0.37) 0.7455
. So from the given options we can say that option (d) is correct.
2 2
Example 5.51 [MT-2009 (2 marks)]: What is the magnitude of the integral 0 (3 x 4 x 2) dx
using single step application of trapezoidal rule?
(a) 9 (b) 16 (c) 18 (d) 36
Solution (b): From the given data, we have a 0 , b 2 and n 1, so we have
h (b a ) n (2 0) 1 2 . So from Eq. 5.36 with f ( x ) 3 x 2 4 x 2 , we have
2
I f ( x ) dx h (1 2) f ( a ) (1 2) f (b) 2 (1 2) f (2) (1 2) f (0) 18 2 16 .
0
1 3
Example 5.53 [AE-2012 (2 marks)]: The integration 0 x dx computed using trapezoidal rule with
n 4 intervals is ……
Solution: From the given data, we have a 0 , b 1 and n 4 , so we have
h (b a ) n (1 0) 4 1 4 . So from Eq. 5.36 with f ( x) x 3 , we have
1
I f ( x) dx h (1 2) f ( a ) f ( x1 ) f ( x2 ) f ( x3 ) (1 2) f (b) , where xi a ih . So
0
3 3 3 3 3 3
1 1 3 1 1 3 1 1 1 1 3 1
I (0) f f f f (1) 0.266 .
4 2 4 2 4 2 4 4 2 4 2
Example 5.54 [CS-2013 (1 mark)]: Function f is known at the following points:
x 0 0.3 0.6 0.9 1.2 1.5 1.8 2.1 2.4 2.7 3.0
f ( x) 0 0.09 0.36 0.81 1.44 2.25 3.24 4.41 5.76 7.29 9.0
3
The value of 0 f ( x )dx computed using the trapezoidal rule is
(a) 8.983 (b) 9.003 (c) 9.017 (d) 9.045
Solution (d): From the given data, we have a 0 , b 3 and we have 10 intervals, i.e., n 10 , and
so we have h (b a ) n (3 0) 10 0.3 . So from Eq. 5.36, we have
3
I f ( x ) dx h (1 2) f (0) f (0.3) f (0.6) f (2.7) (1 2) f (3)
0
I 0.3 (1 2)(0) (0.9 0.36 7.29) (1 2)(9) 0.3 0 (0.9 0.36 7.29) 4.5
I 0.3 0 (25.65) 4.5 9.045 .
[Similar question was also asked in AG-2008 (1 mark)]
Example 5.55 [ME-2014 (2 marks)]: Using the trapezoidal rule, and dividing the interval of
1
integration into three equal subintervals, the definite integral 1 x dx is …………….
Solution: From the given data, we have a 1 , b 1 and n 3 , so we have
h (b a ) n (1 1) 3 2 3 . So from Eq. 5.36 with f ( x) x , we have
1
I f ( x ) dx h (1 2) f (a ) f ( x1 ) f ( x2 ) (1 2) f (b) , where xi a ih . So
1
2 1 2 4 1 2 1 1 1 1
I 1 1 1 1 1.11 .
3 2 3 3 2 3 2 3 3 2
4
Example 5.56 [ME-2014 (2 marks)]: The value of 2.5 ln( x)dx calculated using the Trapezoidal rule
with five sub-intervals is …………….
Solution: From the given data, we have a 2.5 , b 4 and n 5 , so we have
h (b a) n (4 2.5) 5 0.3 . So from Eq. 5.36 with f ( x ) ln x , we have
4
I f ( x )dx h (1 2) f ( a ) f ( x1 ) f ( x2 ) f ( x3 ) f ( x4 ) (1 2) f (b) , where xi a ih . So
2.5
4
I f ( x ) dx 0.3 (1 2) ln(2.5) ln(2.8) ln(3.1) ln(3.4) ln(3.7) (1 2) ln(4) 1.75 .
2.5
3
Example 5.57 [ME-2014 (1 mark)]: The definite integral 1 (1 x) dx is evaluated using Trapezoidal
rule with a step size of 1. The correct answer is …………….
Solution: From the given data, we have a 1 , b 3 and h 1, so we have
h (b a) n 1 (3 1) n n 2 . So from Eq. 5.36 with f ( x) 1 x , we have, x1 a h , so,
3
I f ( x) dx h (1 2) f ( a ) f ( x1 ) (1 2) f (b) 1(1 2)(1) (1 2) (1 2)(1 3) 1.1667 .
1
Example 5.58 [ME-2016 (1 mark)]: Numerical integration using trapezoidal rule gives the best
result for a single variable function, which is
(a) linear (b) parabolic (c) logarithmic (d) hyperbolic
Solution (a): The Trapezoidal rule gives the exact result when applied to any function whose second
derivative is identically zero, i.e., the Trapezoidal rule gives the exact result for polynomials of degree
up to or equal to one, which is linear polynomial. So option (a) is correct.
Example 5.59 [ME-2016 (2 marks)]: The error in numerically computing the integral
0 (sin x cos x)dx using the trapezoidal rule with three intervals of equal length between 0 and is
_____.
Solution: Using trapezoidal rule in [0, ] with three intervals, so n 3 , i xi f ( x)
b
I f ( x) dx , where a0, b ; with step size 0 0 1
a
1 3 1.366
h (b a ) n ( 0) 3 3 , also xi 1 xi h and f ( x ) sin x cos x so
b
2 2 3 0.366
I f ( x) dx h (1 2) f ( a) f ( x1 ) f ( x2 ) (1 2) f (b) , 3 1
a
where x0 a 0 , x5 b .
Thus I f ( x) dx ( 3) (1 2) f (0) f ( 3) f (2 3) (1 2) f ( )
0
I f ( x) dx ( 3) (1 2)(1) 1.366 0.366 (1 2)(1) 1.814 , which is an approximate solution.
0
For exact solution: 0 (sin x cos x)dx ( cos x sin x)0 {(1) 0} {1 0} 2
Example 5.60 [AG-2016 (2 marks)]: Integration by trapezoidal method of log10 ( x ) with lower limit
of 1 to upper limit of 3 using seven distinct values (equally covering the whole range) is _____.
b
Solution: Using trapezoidal rule in [1, 3] , I f ( x) dx , where x0 a 1 , i xi f ( x)
a
0 1.00 0.000
x8 b 3 ; as we have to integrate using seven distinct values (equally covering 1 1.25 0.097
the whole range) so n 7 1 8 , hence the step size is 2 1.50 0.176
h (b a ) n (3 1) 8 0.25 . We have f ( x ) log10 ( x ) , and the values of 3 1.75 0.243
f ( x ) at different points are given in the table thus 4 2.00 0.301
b 5 2.25 0.352
I f ( x) dx h (1 2) f ( x0 ) f ( x1 ) f ( x2 ) f ( x7 ) (1 2) f (b) 6 2.50 0.398
a
3 7 2.75 0.439
I f ( x) dx (0.25) (1 2) f (0) f (1.25) f (2.75) (1 2) f (3.00) 8 3.00 0.477
1
I (0.25) (0.5)(0) 0.097 0.176 0.243 0.301 0.352 0.398 0.439 (0.5)(0.477) 0.5611
A similar formula holds for the next two subintervals from x2 to x4 , and so on. By summing all these
b
m formulas we have, a f ( x) dx ( h 3) f 0 4 f1 2 f 2 4 f3 2 f 4 2 f 2 m 2 4 f 2 m 1 f 2 m
b
or, a f ( x) dx ( h 3) f 0 4 f1 f 3 f 2 m 1 2 f 2 f 4 f 2 m 2 f 2 m (5.43)
where, h (b a) (2m) , f j f ( x j ) . [The formula in Eq. 5.43 was asked in PI-2010 (1 mark)].
Error Estimation of Simpson’s 1/3 Rule: If the fourth derivative f (4) exists and is continuous
on a x b , the error of Eq. 5.43, call it as S1 3 , is given as,
(b a )5 (b a) (2m)
S1 3 4
f (4) (tˆ) h 4 f (4) (tˆ)
h5 f (4) (tˆ) (5.44)
180(2m) 180 180
here tˆ is a suitable unknown value between a and b . With this Eq. 5.43 can be written as,
b h (b a ) 4 (4)
a f ( x)dx 3 f0 4 f1 f3 f 2m 1 2 f2 f 4 f 2m 2 f2 m 180 h f (tˆ) (5.45)
We also note that the error term for the Simpson’s 1/3 rule f (4) , so the Simpson’s 1/3 gives
the exact result when applied to any function whose fourth derivative is identically zero, i.e.,
the Simpson’s 1/3 rule gives the exact result for polynomials of degree up to or equal to three.
[This point was asked in CS-1993, EC-1993, AG-2011 (1 mark)].
The error in Simpson rule is proportional to h5 so it is a third order method, i.e., O ( h5 ) .
[This point was asked in ME-2003 (1 mark)]
Error Bounds in Simpson’s 1/3 Rule: By taking for f (4) in Eq. 5.44, the maximum M 1 and
minimum M 2 on the interval of integration we get the error bounds as (note that C is negative)
(b a )5 (b a )
CM 1 S1 3 CM 2 , where C 4
h4 (5.46)
180(2m) 180
Example 5.62 [AG-2007 (1 mark)]: Integrating the function f ( x ) 1 e x sin(4 x) over the interval
[0,1] using Simpson’s 1/3rd rule gives
(a) 1.021 (b) 1.091 (c) 1.321 (d) 2.642
Solution (c): From the given data we have a 0 , b 1 , and let we have two subintervals so we have
m 1 , h (b a ) 2m (1 0) 2 h 0.5 . So, with f ( x ) 1 e x sin(4 x) , and xi a ih from
1
Eq. 5.43, we have 0 f ( x)dx (h 3) f ( x0 ) 4 f ( x1 ) f ( x2 ) (0.5 3) 1 4(1.551) 0.722 1.321 ,
[Similar questions were also asked in AE-2014 (1 mark), PI-2008, (2 marks)]
Example 5.63 [CH-2008 (2 marks)]: Using Simpson’s 1/3rd rule and FOUR equally spaced intervals
4
( n 4 ), estimate the value of the integral 0 (sin x)
(cos 3 x ) dx
(a) 0.3887 (b) 0.4384 (c) 0.5016 (d) 0.5527
Solution (c): From the given data we have a 0 , b 4 , and we have 4 subintervals, i.e.,
n 2m 4 h (b a) 2m ( 4 0) 4 h 16 . So we have an even number of equal sub-
intervals and hence we apply Simpson’s 1/3 rule. So, with f ( x ) (sin x) (cos3 x ) , from Eq. 5.43, we
4
have I f ( x) dx ( h 3) f ( x0 ) 4 f ( x1 ) f ( x3 ) 2 f ( x2 ) f ( x4 ) , where xi a ih ; thus
0
Statement for Linked Answer Questions 5.64 & 5.65: The following two questions relate to
b
Simpson’s rule for approximating the integral a f ( x) dx on the interval [ a, b] .
Example 5.64 [AE-2008 (2 marks)]: Which of the following gives the correct formula for Simpson’s
rule?
(a) k f (b) f {(a b) 2} (b) k f ( a ) f (b ) 2 f {( a b ) 2}
(c) k f ( a ) f (b ) 3 (4 3) f {( a b ) 2} (d) k f ( a ) f (b ) 3 (4 3) f {( a b ) 3}
where k (b a) 2 .
Solution (c): From the given options we can say that the step size is h (b a ) (2 m ) (b a ) 2
b
m 1 . So from Eq. 5.43, we have a f ( x ) dx ( h 3) f ( x0 ) 4 f ( x1 ) f ( x1 ) , where x0 a ,
ba ab b a f ( a ) f (b) 4 a b
x1 x0 h a and x2 x0 2h b ; thus I f .
2 2 2 3 3 2
Example 5.65 [AE-2008 (2 Solution: Using the result of previous problem we have
marks)]: The percentage error 3
1 0 03 13 4 1
(with respect to the exact solution) I 0.25 . Also the exact value of the
in estimation of the integral 2 3 32
1 1
0 x 3 dx using Simpson’s rule is given integral is I E x3 dx ( x 4 4)10 1 4 0.25 . So the
0
(a) 5.3 (b) 3.5 percentage error with respect to the exact solution is
(c) 2.8 (d) 0 ( I E I ) I E 100 0 .
1 2
Example 5.66 [XE-2009 (2 marks)]: Simpson’s 1/3 rule applied to 1 (3 x 5)dx , with sub-interval
h 1, will give
(a) the exact result (b) error between 0.01% to 0.1%
(c) error between 0.1% to 1.0% (d) error 1.0%
2
Solution: (a): As the given function f ( x ) 3x 5 is a 2nd degree polynomial and we know that
Simpson’s 1/3 rule gives the exact result for polynomials of degree up to or equal to three. So option
(a) is correct.
Example 5.67 [CE-2010 (2 marks)]: The table below x 0 0.25 0.50 0.75 1.0
gives values of a function f ( x ) obtained for values of x f ( x ) 1 0.9412 0.8 0.64 0.50
at intervals of 0.25.
The value of the integral of the function between the limits 0 to 1 using Simpson’s rule is
(a) 0.7854 (b) 2.3562 (c) 3.1416 (d) 7.5000
Solution (a): From the given data we have a0, b 1, h 0.25
h (b a ) 2m 0.25 (1 0) 2m m 2 . So we have an even number of equal sub-intervals
and hence we apply Simpson’s 1/3 rule. Thus from Eq. 5.43, we have
1
0 f ( x)dx (h 3) f ( x0 ) 4 f ( x1 ) f ( x3 ) 2 f ( x2 ) f ( x4 ) , where xi a ih ; So from the given
1
data we have f ( x) dx (0.25 3) 1 4 0.9412 0.64 2(0.8) 0.50 0.7854
0
[Similar questions were also asked in TF-2007, MN-2007, TF-2012 (2 marks)]
1.5
Example 5.68 [CE-2012 (1 mark)]: The estimate of 0.5 (1 x)dx obtained using Simpson’s rule with
three-point function evaluation exceeds the exact value by
(a) 0.235 (b) 0.068 (c) 0.024 (d) 0.012
Solution (d): As three-point function evaluation means we have an even number of equal sub-
intervals, i.e., with m 1 , we have n 2 , and so we apply Simpson’s 1/3 rule. Thus with a 0.5
and b 1.5 , h (b a ) 2m (1.5 0.5) 2 0.5 . So, with f ( x ) 1 x , from Eq. 5.43, we have
1.5
0.5 f ( x)dx (h 3) f ( x0 ) 4 f ( x1 ) f ( x2 ) ; as x0 a 0.5 , x1 a h 1 and x2 b 1.5 . So
1.5 0.5 1 1 1 1.5 1 1.5
0.5 f ( x)dx
3 0.5
4
1 1.5
1.111 . As 0.5 x dx (ln x)0.5 ln1.5 ln 0.5 1.099 , which is
1.5 1.5
the exact value of 0.5 (1 x)dx . So the estimate of 0.5 (1 x)dx obtained using Simpson’s rule with
three-point function evaluation exceeds the exact value by (1.111 1.099) 0.012 .
[Similar questions were also asked in EE-1998, ME-2011, MN-2009 (2 marks)]
4 2
Example 5.69 [XE-2012 (2 marks)]: The exact solution of the integral 0 ( x 4)dx is denoted by
I E . The same integral evaluated numerically by the trapezoidal rule and the Simpson’s 1/3 rule are
denoted by IT and I S , respectively. The subinterval used in the numerical method is h 2 , then
(a) I E I S IT (b) I E I S IT (c) I E I S IT (d) IT I S I E
4 4
Solution (b): I E ( x 2 4) dx ( x 3 3) 4 x 16 3 . For IT , we have a 0 , b 4 and h 2
0 0
2 (4 0) n n 2 ;
IT h (1 2) f ( a) f ( a h) (1 2) f (b) 2 (1 2)( 4) 0 (1 2)(12) 8 . For I S , as we have 2nd
degree polynomial and we know that Simpson’s 1/3 rule gives the exact result for polynomials of
degree up to or equal to three. So I S 16 3 . Thus we have I E I S IT .
Example 5.70 [CE-2013 (2 marks)]: The magnitude as the error (correct to two decimal places) in
4 4
the estimation of integral 0 ( x 10) dx using Simpson 1 3 rule, by taking the step length as 1, is ___
Solution: From the given data we have a0, b 4, h 1
h (b a) 2m 1 (4 0) 2m m 2 . So we have an even number of equal sub-intervals and
hence we apply Simpson’s 1 3 rule. So, with f ( x ) x 4 10 , from Eq. 5.43, we have
4
I S f ( x )dx ( h 3) f ( x0 ) 4 f ( x1 ) f ( x3 ) 2 f ( x2 ) f ( x4 ) , where xi a ih ; thus
0
I S (1 3) f (0) 4 f (1) f (3) 2 f (2) f (4) (1 3) 10 4 11 91 2(26) 266 245.33 .
4
Now the exact value of the given integral is I E ( x 4 10) dx {( x 5 5) 10 x}40 244.8 . So the
0
magnitude of error in the estimation of the given integral using Simpson’s 1 3 rule is
I E I S 244.8 245.33 0.53 .
0.5 x3
Example 5.71 [CH-2013 (2 marks)]: The value of the integral 0.1 e dx evaluated by Simpson’s
rule using 4 subintervals (up to 3 digits after the decimal point is) ……
Solution: From the given data we have a 0.1 , b 0.5 , and we have 4 subintervals, i.e., n 2m 4
h (b a) 2m (0.5 0.1) 4 h 0.1 . So we have an even number of equal sub-intervals and
3
hence we apply Simpson’s 1/3 rule. So, with f ( x ) e x , from Eq. 5.43, we have
2
1 f ( x)dx (h 3) f ( x0 ) 4 f ( x1 ) f ( x3 ) 2 f ( x2 ) f ( x4 ) , where xi a ih ; thus
0.5
0.1 f ( x)dx (0.1 3) 0.999 4 0.992 0.938 2(0.973) 0.882 0.385 .
[Similar question was also asked in AG-2013 (2 marks)]
Example 5.72 [CS-2014 (2 marks)]: With respect to the numerical evaluation of the definite integral,
b
K x 2 dx , where a and b are given, which of the following statements is/are TRUE?
a
A. The value of K obtained using the trapezoidal rule is always greater than or equal to the exact
value of the definite integral
B. The value of K obtained using the Simpson’s rule is always equal to the exact value of the
definite integral
(a) A only (b) B only (c) Both A and B (d) Neither A or B
Solution (c): Let I E and IT be the exact solution and evaluated by Trapezoidal rule, respectively, of
the given integral. As from Eq. 5.39, we can write IT I E (b a ) 12 h 2 f (tˆ) ; so we can say that
I T I E , where equality holds when the given function is of at most degree one. Also the given
function is 2nd degree polynomial and we know that Simpson’s 1/3 rule gives the exact result for
polynomials of degree up to or equal to three. So I S I E , where I S is the integral evaluated by using
Simpson’s rule. So both statements are correct.
Example 5.73 [MA-2014 (2 marks)]: Let the x: 0 0.25 0.5 0.75 1.0
following discrete data be obtained from a curve y : 1 0.9896 0.9589 0.9089 0.8415
y y ( x ) . Let S be the solid of revolution obtained by rotating the above curve about the x axis
between x 0 and x 1 and let V denote its volume. The approximate value of V , obtained using
Simpson’s 1/3rd rule, is …………….
Solution: If a curve y f ( x) , between x 0 and x 1 , is rotated about x axis then the volume
1 1
generated will be V y 2 dx y 2 dx . Now from the given data we have a 0 , b 1 , we have
0 0
Example 5.74 [PI-2014 (2 marks)]: Using the Simpson’s 1/3rd rule, the x 1 3 5
y 2 6 4
value of ydx computed, for the data given below is ………
Solution: From the given data we have a 1 , b 5 , and let we have two subintervals so we have
m 1 , h (b a) 2m (5 1) 2 h 2 . So, with y f ( x) , from Eq. 5.43, we have
1
0 ydx (h 3) f (1) 4 f (3) f (5) (2 3) 2 4(6) 4 20 .
2 x
Example 5.75 [TF-2014 (2 marks)]: 0 e dx is evaluated both by trapezoidal rule and Simpson’s
1/3rd rule by taking two subintervals. The difference between the results, accurate to the second
decimal place, is …………….
Solution: From the given data we have a 0 , b 2 , and n 2 h (b a) n (2 0) 2 1 . So
evaluation of the given integral by Trapezoidal rule is
2
IT e x dx h (1 2) f ( a ) f ( x1 ) (1 2) f (b) 1 (1 2) e (1 2) e 2 6.913 , where x1 a h .
0
Now for Simpson’s 1/3 rule, we have n 2m m 1 , so
2
I S e x dx ( h 3) f ( a ) 4 f ( x1 ) f (b ) (1 3) 1 4e e 2 6.420 , where x1 a h . So
0
A similar formula holds for the next two subintervals from x3 to x6 , and so on. By summing all these
m formulas we obtain Simpson’s rule as,
I (3h 8) f 0 3 f1 3 f 2 f 3 f 3 3 f 4 3 f5 f 6 ..... f n 3 3 f n 2 3 f n 1 f n (5.49)
n 2 n 1 n 3
I (3h 8) f 0 3 i 1,4,7,... fi 3i 2,5,8,... fi 2 i 3,6,9,... fi f ( xn ) (5.50)
where, h (b a ) (3m) and f j f ( x j ) .
Error Estimation of Simpson’s 3/8 Rule: If the fourth derivative f (4) exists and is continuous
on a x b , the error of Eq. 5.50, call it as S3 8 , is given as,
Forward Euler’s Method: Forward Euler’s method (or simply Euler method) provides an
approximation of the solution to a first order ODE given by Eq. 5.57. From a graphical perspective, a
solution of first-order differential equation is a function whose graph is consistent with the direction
field. In other words, the graph of the solution is tangent to the slope lines at every point. Euler’s
method approximates the solution by a sequence of points which follows the slope lines. Given Eq.
5.54 and the initial condition y ( x0 ) y0 , Euler’s method approximates the solution at x1 , x2 , ,
where the consecutive x ’s differ by h , so that xi x0 ih , i 1, 2, . The increment h is called step
size, which is a parameter of the method, which determines the accuracy of the approximation. Let yi
denote the approximated value of the solution at xi . The idea is to approximate the solution at each
step by the slope line. Consider Fig. 5.10, we start at the
initial point ( x0 , y0 ) . The straight line segment is the
tangent line to the solution curve at ( x0 , y0 ) . Its slope is
y ( x0 ) , which is given by Eq. 5.54, i.e., y ( x0 ) f ( x0 , y0 )
. Hence Forward Euler method suggests that in Eq. 5.55,
2 3
for a small h the higher powers h , h , are very small
which suggests the approximation
y ( x h) y ( x ) hy ( x) y ( x ) hf ( x, y ) and the
st
following iteration process. In the 1 step we compute
y1 y0 h ( dy dx ) ( x , y ) y0 hf ( x0 , y0 ) which
0 0
in general we have,
yn 1 yn h (dy dx ) ( x yn h f ( xn , yn ) , ( n 0,1, ) (5.56)
n , yn )
[Eq. 5.56 was asked in AE-2007 (1 mark)]. This method is hardly ever used in practice, but since it
is simple, it is nicely explains then principle of methods based on the Taylor series. In Euler’s method,
the sequence of lines may also deviates considerably from the curve of solution. Hence there is a
modification of this method; two such methods that are of practical importance, namely, ‘improved
Euler method’ and ‘Runge – Kutte method’ which are discussed later.
Example 5.78 [AG-2007 (1 mark)]: Taking y (0) 0 and using Euler’s method with step size
h 0.1 solution of the differential equation dy dx 2 xy 1 gives the value of y (0.3) as
(a) 0.3101 (b) 0.3142 (c) 0.6202 (d) 4.0800
Solution (a): From the given data we have f ( xn , yn ) ( dy dx ) ( x , y ) 2 xn yn 1 , h 0.1 , y (0) 0 .
n n
Example 5.79 [IN-2010 (2 marks)]: The velocity v (in m/s) of a moving mass, starting from rest, is
given as dv dt v t . Using Euler forward difference method (also known as Cauchy-Euler method)
with a step size of 0.1 sec, the velocity at 0.2 sec evaluates to
(a) 0.01 m/s (b) 0.1 m/s (c) 0.2 m/s (d) 1 m/s
Solution (a): From the given data we have f (tn , vn ) ( dv dt ) ( t , v ) vn tn , h 0.1 , v (0) 0 . So
n n
Example 5.80 [XE-2013 (2 marks)]: Using Euler’s method to solve the differential equation
dy dx 2 cos(4 x 3) y , y (0) 1 with step-size h 0.25 , the value of y (0.5) is
(a) 1.3125 (b) 1.1875 (c) 1.125 (d) 1.0625
Example 5.81 [CH-2014 (2 marks)]: Consider the following differential equation dy dx x ln( y ) ;
y 2 at x 0 . The solution of this equation at x 0.4 using Euler method with a step size of
h 0.2 is …………….
Solution: From the given data we have f ( xn , yn ) ( dy dx ) ( x , y ) xn ln( yn ) , h 0.2 , y (0) 2 .
n n
Example 5.82 [MA-2014 (1 mark)]: Let y ( x) be the solution to the initial value problem
dy dx y 2 x subject to y (1.2) 2 . Using the Euler method with the step size h 0.05 , the
approximate value of y (1.3) , correct to two decimal places, is …………….
Solution: From the given data we have f ( xn , yn ) ( dy dx) ( x yn 2 xn , h 0.05 , y (1.2) 2
n , yn )
[Eq. 5.57 was asked in CS-1994 (1 mark)]. Normally we do not know the derivative at point x h ,
although we need it to compute the function value at point x h . In practice this requires a
rearrangement of the equation. We call such a numerical scheme an implicit numerical scheme, as
opposed to explicit numerical schemes, such as the Forward Euler method, where no such
rearrangement is necessary.
Example 5.83 [CE-2006 (1 mark)]: The differential equation dy dx 0.25 y 2 is to be solved using
the backward (implicit) Euler’s method with the boundary condition y 1 at x 0 and with a step
size of 1. What would be the value of y at x 1 ?
(a) 1.33 (b) 1.67 (c) 2.00 (d) 2.33
2
Solution: From the given data we have f ( xn , yn ) ( dy dx ) ( x , y ) 0.25 yn , h 1, y (0) 1 . So from
n n
Improved Euler’s Method: In the improved Euler method or improved Euler – Cauchy method,
in each step we compute first the auxiliary value, i.e.
yn*1 yn hf ( xn , yn ) (5.58)
and then the new value is given as,
yn 1 yn (1 2) h f ( xn , yn ) f ( xn 1 , yn*1 ) (5.59)
This method has a simple geometric interpretation. In fact, we may say that in the interval from xn to
xn (1 2) h we approximate the solution y by the straight line through ( xn , yn ) with slope f ( xn , yn )
, and then we continue along the straight line with slope f ( xn 1 , yn* 1 ) until x reaches xn 1 . The
improved Euler – Cauchy method is predictor – corrector method, because in each step we first
predict a value by Eq. 5.58 and then correct it by Eq. 5.59.
Example 5.84 [IN-2013 (2 marks)]: While numerically solving the differential equation
( dy dx ) 2 xy 2 0 , y (0) 1 using Euler’s predictor-corrector (improved Euler-Cauchy) method
with a step size of 0.2, the value of y after the first step is
(a) 1.00 (b) 1.03 (c) 0.97 (d) 0.96
2
Solution (d): From the given data, we have h 0.2 , y (0) 1 , f ( x, y ) (dy dx) ( x, y ) 2 xy ; so
2
f ( xn , yn ) 2 xn yn2 and
*
*
f ( xn 1 , yn 1 ) 2 xn 1 yn 1 , where from Eq. 5.58, we have
yn*1 yn hf ( xn , yn ) yn h ( 2 xn yn2 ) yn 2hxn yn2 ; f ( xn 1 , yn* 1 ) 2 xn 1 ( yn 2 hxn yn2 ) 2 . So
from Eq. 5.59, we have yn 1 yn (1 2)h 2 xn yn2 2 xn 1 ( yn 2hxn yn2 ) 2
, where xi x0 ih . Thus
starting with x0 0 and y ( x0 ) 1 , we have y1 y0 (1 2) h 2 x0 y02 2 x1 ( y0 2hx0 y02 ) 2
y1 1 (1 2)0.2 0 2(0.2)(1 0) 2 1 (1 2 )0.2( 0.4) 1 0.04 0.96 .
Runge-Kutta (RK) Method: The Taylor’s series method of solving differential equations
numerically is restricted by the finding the higher order derivatives. However there is a class of
methods known as Runge – Kutta methods which do not require the calculations of higher order
derivatives. These methods agree with Taylor’s series solution up to the terms in h r , where r differs
from method to method and is called the order of that method. The Euler’s method and improved
Euler’s method are the Runge – Kutta methods of first and second order, respectively. The Runge –
Kutta methods of order three are not generally used. A method of great practical importance and much
greater accuracy than that of improved Euler’s method is the classical Runge – Kutta method of fourth
order. For find the solution of y dy dx f ( x, y ) , y ( x0 ) y0 at equidistant points
x1 x0 h, x2 x0 2h, , x N x0 Nh . We see that in each step we first compute four auxiliary
quantities k1 , k2 , k3 , k 4 and then the new value yn 1 at xn 1 is given as,
Example 5.86 [MN-2011 (2 marks)]: For the equation dy dx 2 x 3 y , the value of y at x 0.1
in one step using Runge-Kutta fourth order method for the condition y 1 when x 0 , is
(a) 0.3608 (b) 1.2508 (c) 1.3608 (d) 1.4625
Solution (c): From the given data we have, f ( x, y ) 2 x 3 y , y 1 at x 0 and h 0.1 . So from
Eq. 5.60, we have y1 y0 (1 6)( k1 2k2 2k3 k4 ) , where x1 x0 h , k1 hf ( x0 , y0 ) ,
k2 hf x0 (1 2)h, y0 (1 2) k1 , k3 hf x0 (1 2)h, y0 (1 2) k2 and k4 hf ( x0 h, y0 k3 ) . So
k1 0.1(2 x0 3 y0 ) 0.1(2 0 3 1) 0.3 ;
k2 0.12 x0 (1 2) h 3 y0 (1 2) k1 0.10.1 3 1 (0.3 2) 0.355 ;
k3 0.12 x0 (1 2) h 3 y0 (1 2) k2 0.10.1 3 1 (0.355 2) 0.363 ;
k4 0.1{2( x0 h) 3( y0 k3 )} 0.1{0.2 3(1 0.363)} 0.429 . So
y1 1 (1 6)(0.3 2(0.355) 2(0.363) 0.429) 1.360 .
[Similar questions were also asked in CS-1993 (1 mark), XE-2008, MA-2014 (2 marks)]
Example 5.87 [XE-2007 (2 marks)]: If Runge-Kutta method of order 4 is used to solve the
differential equation dy dx f ( x ) , y (0) 0 in the interval [0, h] with step size h , then
(a) y ( h) ( h 6) f (0) 4 f ( h 2) f ( h) (b) y ( h) ( h 6) f (0) f ( h)
(c) y ( h) ( h 2) f (0) f (h) (d) y ( h) ( h 6) f (0) 2 f ( h 2) f ( h)
Solution (a): To avoid confusion with the function f , let us consider it is g and at the end we will
replace it by f . So we can say that dy dx f ( x ) ; thus we have, f ( x, y ) g ( x ) , y 0 at x 0 and
step size h . So from Eq. 5.60, we have y1 y0 (1 6)( k1 2k2 2k3 k4 ) , where x1 x0 h ,
k1 hf ( x0 , y0 ) , k2 hf x0 (1 2)h, y0 (1 2) k1 , k3 hf x0 (1 2)h, y0 (1 2) k2 and
k4 hf ( x0 h, y0 k3 ) . So k1 h g ( x0 ) h g (0) ; k2 h g x0 (h 2) h g ( h 2) ;
k3 h g x0 ( h 2) h g ( h 2) ; and k4 h g ( x0 h) h g ( h) . Thus
y1 y0 (h 6) g (0) 2 g ( h 2) 2 g (h 2) h g ( h) 0 ( h 6) g (0) 4 g ( h 2) g (h) . So replacing
g by f , we get y1 ( h 6) f (0) 4 f ( h 2) f ( h) .
Error in Numerical Methods: The numerical methods are simply approximations of the true
solution, and as such they inevitably introduce errors. We can classify the errors into three types: (1)
The local error describes the error that occurs after a single time step of the method. (2) The global
error describes the cumulative effect of the local errors after many time steps. (3) The roundoff error
arises from the finite precision with which computers store and process numbers. Let the set
{( xn , yn ) : n 0,1,} be an approximation to the solution y ( x ) of the initial value problem. The
global error at xn is the difference between the true solution at xn and the approximation yn , i.e.,
En y ( xn ) yn . The absolute error in the numerical approximation at x xn is En , i.e., the
absolute value of the global error at xn . The local error is the difference between the true solution
y ( xn 1 ) and the approximation after one time step, assuming that ( xn , yn ) is a point on the graph of
the true solution (i.e. assuming that yn is exact). A numerical method is of order n , where n is a
positive integer, if the method is exact for polynomials of degree n or less. In other words, if the true
solution of an initial-value problem is a polynomial of degree n or less, then the approximate solution
and the true solution will be identical for a method of order n . In general, higher the order, the more
accurate will be the method.
Error in Forward Euler Method: Consider the Forward Euler method. The local error at the nth
step is given by Enloc y ( xn 1 ) { yn h f ( xn , yn )} . As from Taylor’s theorem, we have
y ( xn h ) y ( xn ) hy ( xn ) (1 2) h 2 y ( ) yn h f ( xn , yn ) (1 2) h 2 y ( ) , where is the value
between x0 and x0 h . Hence the local error in Forward Euler’s method is E loc (1 2) h 2 y ( ) h 2
[This point was asked in PI-2009 (1 mark)]. Hence the local error in Forward Euler’s method
decreases as the square of the step size. Thus if the step size is cut in half, then the local error at each
step is reduced by a factor of four. A derivation of the global error is beyond the scope of this text. We
shall restrict ourselves to a simple argument. Suppose that y (0) y0 and we want to find y ( X ) . The
number of steps required is approximately X h . If the local error is roughly constant for 0 x X ,
and E( X ) is the global error at x X then
E ( X ) number of iterations E loc E ( X ) ( X h) h 2 E ( x ) h . So the global error in
Forward Euler’s method decreases linearly with the step size. Thus, we can say that Forward Euler’s
method is a first-order method, because the global error is proportional to h .
Example 5.88 [PI-2010 (2 mark)]: Euler’s method of integration is applied to the initial value
problem: dy dx 2 x ; y (0) 0 . If the step size h 0.2 , then the error in computation (in
percentage) after 5 steps would be
(a) 0 (b) 10 (c) 20 (d) 30
Solution (c): From the given data, y 2 x y 2 , y (0) 0 , i.e., x0 0 , h 0.2 . So the local
error is E loc (1 2)h 2 y( ) (1 2) 0.22 2 0.04 . If n is the number of iterations then, global
error at x X x0 5h 0 5(0.2) 1 is E global n E loc ( X h) E loc (1 0.2) 0.04 20% .
Error for Backward Euler’s method: To determine the local error for the Backward Euler’s
method we expand y( x) in a Taylor series about xn 1 , i.e.,
y ( xn 1 h ) y ( xn 1 ) hy( xn 1 ) (1 2) h 2 y ( x ) . But y ( xn 1 h) y ( xn ) yn ; y ( xn 1 ) y ( xn h)
and y ( xn 1 ) f ( xn 1 , yn 1 ) . Thus we have y ( xn h ) yn h f ( xn 1 , yn 1 ) (1 2) h 2 y ( x ) . As
Enloc y ( xn 1 ) { yn h f ( xn 1 , yn 1 )} Enloc (1 2) h 2 y ( x ) h 2 . Hence local error in Backward
Euler’s method decreases as the square of the step size. Now with same argument, as with Forward
Euler’s method, the global error in Backward Euler’s method decreases linearly with the step size.
Error of the improved Euler method: The local error is of order h3 and the global error is of
order h 2 , so that the method is a second – order method.
Error in Runge-Kutta method: This method has local truncation error O ( h 4 ) . The method is
well suited to the computer because it needs no special starting procedure, makes light demand on
storage, and repeatedly uses the same straightforward computational procedure. It is numerically
stable. It is to be also noted that if f depends only on x , this method reduces to Simpson’s rule of
integration. Note further that k1 , , k4 depends on n and generally change from step to step.
Example 5.89 [EC-2016 (1 mark)]: Consider the first order initial value problem y y 2 x x 2 ,
y (0) 1 , (0 x ) , with exact solution y ( x ) x 2 e x . For x 0.1 , the percentage difference
between the exact solution and the solution obtained using a single iteration of the second-order
Runge-Kutta method with step-size h 0.1 is _____.
Solution: For exact solution: y (0.1) (0.1) 2 e0.1 1.1152
For solution using second-order Runge-Kutta method, which is improved Euler’s method, with step-
size h 0.1 , with initial guess x0 0 and y0 1 , is given by
y1 y0 (1 2) h f ( x0 , y0 ) f ( x0 h, y1* ) , where y1* y0 hf ( x0 , y0 ) and f ( x, y ) y 2 x x 2 .
As y1* y0 hf ( x0 , y0 ) 1 (0.1) f (0,1) 1 0.1{1 2(0) 0 2 } 1.1 , so
y1 1 (1 2)(0.1) f (0,1) f (0.1,1.1) 1 (1 2)(0.1) {1 2(0) 02 } {1.1 2(0.1) 0.12 } 1.1145
yExact 1.1152 , yApprox. 1.1145 ,
Exact Approx. 1.1152 1.1145
% Error 100 100 0.063%
Exact 1.1152
Example 5.90 [TF-2016 (1 mark)]: Which of the following is a multi-step numerical method for
solving the ordinary differential equation?
(a) Euler method (b) Improved Euler method
(c) Runge-Kutta method (d) Adams-Multon method
Solution (d): The Euler method, Improved Euler method and Runge-Kutta method are not multi-step
numerical method for solving the ODEs. So, option (d) is correct.
Stability Analysis: If the effect of round – off error remains bounded as n , with a fixed step
size, then the method is said to be stable, otherwise unstable. Unstable method will diverge away from
solution and causes overflow error. Using a general single step method equation, yn 1 E y , then
condition for absolute stability is E 1 , using a test equation y y . Let us find the condition for
stability for Euler’s method, yn 1 yn hf ( xn , yn ) yn 1 yn h yn (1 h ) yn . Now comparing
it with the condition, E 1 , we get, 1 h 1 1 1 h 1 2 h 0 , which is the
required condition for stability of Euler’s method.
Example 5.91 [IN-2009 (2 marks)]: The differential equation dx dt (4 x ) , with x(0) 0 , and
the constant 0 , is to be numerically integrated using the forward Euler method with a constant
integration time step T . The maximum value of T such that the numerical solution of x converges is
(a) 4 (b) 2 (c) (d) 2
Solution: From the given data we have f (tn , xn ) ( dx dt ) (t (4 xn ) , x(0) 0 , and step size
n , xn )
|
an1 an 2 ann xn bm am1 am 2
| bm amn
A solution of Eq. 5.61 is a set of number x1 , x2 , , xn that satisfy all the n equations, and a solution
vector of Eq. 5.62 is a vector x whose components are the solution of Eq. 5.61. The method of
solving such a system by determinants is not metrical, even with efficient methods for evaluating the
determinants.
Gauss Elimination Method: A practical method for the solution of linear system is the so called
Gauss Elimination method. This standard method, for solving linear system (Eq. 5.61), was already
discussed in Linear Algebra chapter. We reduce a given linear system (Eq. 5.61) to triangular form. In
the first step we eliminate x1 from equations E2 to En in Eq. 5.61. We do this by adding (or
subtracting) suitable multiples of E1 from equations E2 , , En and taking the resulting equations,
call them E2* , , En* as the new equations. The first equation E1 is called the pivot equation in this
step, and a11 is called the pivot. This equation is left unaltered. In the second step we take the new
second E 2* (which no longer contains x1 ) as the pivot equation and use it to eliminate x2 from E3* to
E n* . And so on. After n 1 steps this gives a triangular system that can be solved by back
substitution. The pivot akk (in step k ) must be different from zero and should be large in absolute
value, to avoid round off magnification by the multiplication in the elimination. For this we choose as
our pivot equation one that has absolutely largest a jk in column k on or below the main diagonal
(actually, the uppermost of there are several such equations). This popular method is called partial
pivoting. Partial pivoting distinguished it from total pivoting, which involves both row and column
interchanges but is hardly used in practice.
E1 : 8 x2 2 x3 7
Example 5.93: Solve the linear system of equations, E2 : 3 x1 5 x2 2 x3 8
E3 : 6 x1 2 x2 8 x3 26
Solution: The given linear system of equations can be written in form of augmented matrix as,
E1 : 0 8 2 | 7
E2 : 3 5 2 | 8 . We must pivot since E has no x term. In the first column, equation E
1 1 3
E3 : 6 2 8 | 26
has the largest coefficient. Hence we interchange E1 and E3 as,
E3 : 6 x1 2 x2 8 x3 26 E3 : 6 2 8 | 26
E2 : 3 x1 5 x2 2 x3 8 or E2 : 3 5 2 | 8
E1 : 8 x2 2 x3 7 E1 : 0 8 2 | 7
Hence the first pivot equation is E3 and pivot element is 6.
To eliminate x1 from E2 (as x1 is already eliminated from E1 ), we have E2 E2 (1 2) E3 as,
E3 : 6 x1 2 x2 8 x3 26 E3 : 6 2 8 | 26
E2 : 4 x2 2 x3 5 or E2 : 0 4 2 | 5
E1 : 8 x2 2 x3 7 E1 : 0 8 2 | 7
Now, the largest coefficient in the second column is 8, hence we interchange E2 and E1 as,
E3 : 6 x1 2 x2 8 x3 26 E3 : 6 2 8 | 26
E1 : 8 x2 2 x3 7 or E1 : 0 8 2 | 7
E2 : 4 x2 2 x3 5 E2 : 0 4 2 | 5
Hence the second pivot equation is E1 and pivot element is 8.
To eliminate x2 from E2 , we have E2 E2 (1 2) E1 as,
E3 : 6 x1 2 x2 8 x3 26 E3 : 6 2 8 | 26
E1 : 8 x2 2 x3 7 or E1 : 0 8 2 | 7
E2 : 3x3 3 2 E2 : 0 0 3 | 3 2
The resulting triangular system is shown above and this is the end of forward elimination. Now for
back substitution from this system, king the last equation, then second equation, and finally the first
1 1 1
equation, we compute the solution as, x3 ; x2 (7 2 x3 ) 1; x1 (26 2 x2 8 x3 ) 4
2 8 6
Number of operation count in Gauss Elimination method: The important factors in judging
the quality of a numeric method are (1) amount of storage (2) amount of time which is equivalent to
number of operations (3) effect of round off error. For the gauss elimination, the operation count for a
full matrix (a matrix with relatively many non-zero entries) is as follows. In step k we eliminate xk
from n k equations. This needs n k divisions in computing the m jk and ( n k )(n k 1)
multiplications and as many subtractions. Since we do n 1 (as k goes from 1 to n 1 ) and thus the
total number of operations in this forward elimination is given as,
n 1 n 1 n 1 n 1
1 2 2
f ( n) ( n k ) 2 ( n k )( n k 1) s 2 s ( s 1) ( n 1) n ( n 2 1) n n 3 , where
k 1 k 1 s 1 s 1 2 3 3
2 n3 3 is obtained by dropping lower powers of n . We see that f (n) grows about proportional to n3
We say that f (n) is of order n3 and write f ( n) O( n3 ) , where O suggests order. In the back
substitution of xi , we make n i multiplications and as many subtraction, as well as 1 division.
It can be proved that for any non-singular matrix, the rows can be reordered so that the resulting
matrix A has an LU factorization (Eq. 5.63) in which L turns out to be the matrix of the
multipliers m jk of the Gauss elimination, with main diagonal 1,1, ,1 and U is the matrix of the
triangular system at the end of the Gauss elimination. The crucial idea now is that L and U in Eq.
5.63 can be computed directly, without solving simultaneous equations (thus, without using the Gauss
elimination). As a count shows, this need about n3 3 operations, about half as many as the Gauss
elimination, which needs about 2n3 3 . And once we have Eq. 5.63, we can use it for solving Ax b
in two steps, involving only about n 2 operations, simply by noting that Ax LUx b is written as,
Ly b , where Ux y (5.64)
and solving first equation of Eq. 5.64 for y and then second equation of Eq. 5.64 for x . Here we can
require that L have main diagonal 1,1, ,1 , then this is called Doolittle’s method. Both the equations
of Eq. 5.64 are triangular, so we can solve them as in the back substitution for the Gauss elimination.
A similar method, Crout’s method is obtained from Eq. 5.63 if U (instead of L ) is required to have
main diagonal 1,1, ,1 . In either case the factorization is unique. For the solution of Eq. 5.61, the
entries of the matrix L [l jk ] (with main diagonal 1,1, ,1 and l jk suggesting multiplier) and
U [u jk ] in the Doolittle method are computed as,
u1k a1k k 1, 2, , n
l j1 a j1 u11 j 2, 3, , n
j 1 (5.65)
u jk a jk s 1 l js usk k j , , n; j 2
k 1
l jk (1 ukk ) a jk s 1 l js usk j k 1, n; k 2
2 1
Example 5.94 [EE-2011 (2 marks)]: The matrix A is decomposed into a product of a
4 1
lower triangular matrix L and an upper triangular matrix U . The properly decomposed L and U
matrices respectively are
1 0 1 1 2 0 1 1
(a) and (b) and
4 1 0 2 4 1 0 1
1 0 2 1 2
1 0.5 0
(c) and (d)
and
4 1 0 1 3 4
0 1
l11 0 1 u12 2 1
Solution (d): From LU factorization method, we have .
l21 l22 0 1 4 1
l11 2 , l21 4 , l11u12 1 u12 1 2 ; l21u12 l22 1 l22 ( 1 2) 3 . Hence the matrix
2 0 1 0.5 2 0 1 0.5
A LU A L and U .
4 3 0 1 4 3 0 1
[Similar question was also asked in EE-1997 (5 marks)]
Solution of Linear System of Equations by Iterations: The methods discussed in the previous
section belong to the direct methods for solving systems of linear equations; these are methods that
yield solutions after an amount of computations that can be specified in advance. In this section, we
discuss indirect or iterative methods in which we start from an initial value and obtain better and
better approximations from a computational cycle repeated as often as may be necessary, for
achieving a required accuracy, so that the amount of arithmetic depends upon the accuracy required.
Consider a linear system of n linear equations in n unknowns, given by Eq.5.61 in which the
diagonal elements aii do not vanish. Now the system given by Eq. 5.61 can be written as
x1 (b1 a11 ) ( a12 a11 ) x2 ( a13 a11 ) x3 ( a1n a11 ) xn
x2 (b2 a22 ) ( a21 a22 ) x1 ( a23 a22 ) x3 ( a2 n a22 ) xn
x3 (b2 a22 ) ( a21 a22 ) x1 ( a23 a22 ) x2 ( a2 n a22 ) xn (5.66)
xn (bn ann ) (an1 ann ) x1 ( an 2 ann ) x2 ( an ,n 1 ann ) xn 1
Suppose we start with x1(0) , x2(0) , , xn(0) as initial values to the variables x1 , x2 , , xn . Then we can
find better approximations to x1 , x2 , , xn using the following two iterative methods: (a) Jacobi’s
Iteration method (b) Gauss Seidel Iteration method (c) Newton Raphson Method.
Jacobi’s Iteration method: The steps for Jacobi’s Iteration method are given as:
Step 1: Determination of first approximation x1(1) , x2(1) , , xn(1) using x1(0) , x2(0) , , xn(0) and
(1) (0) (0) (0)
x1 (b1 a11 ) ( a12 a11 ) x2 ( a13 a11 ) x3 (a1n a11 ) xn
x2(1) (b2 a22 ) ( a21 a22 ) x1(0) ( a23 a22 ) x3(0) ( a2 n a22 ) xn(0)
(1) (0) (0) (0)
x3 (b2 a22 ) ( a21 a22 ) x1 ( a23 a22 ) x2 ( a2 n a22 ) xn (5.67)
xn(1) (bn ann ) (an1 ann ) x1(0) ( an 2 ann ) x2(0) ( an ,n 1 ann ) xn(0)
1
Step 2: Similarly, x1(2) , x2(2) , , xn(2) are evaluated by just replacing xr(0) (1 r n) in Eq. 5.67 by
xr(1) (1 r n) . Similarly,
Step ( + ): In general, if x1( n ) , x2( n ) , , xn( n ) are a system of nth approximations, then the next
approximation is given by the formula
x1( n 1) (b1 a11 ) ( a12 a11 ) x2( n ) ( a13 a11 ) x3( n ) ( a1n a11 ) xn( n )
x2( n 1) (b2 a22 ) (a21 a22 ) x1( n ) ( a23 a22 ) x3( n ) ( a2 n a22 ) xn( n )
x3( n 1) (b2 a22 ) (a21 a22 ) x1( n ) ( a23 a22 ) x2( n ) ( a2 n a22 ) xn( n ) (5.68)
( n 1) (n) (n) (n)
xn (bn ann ) ( an1 ann ) x1 ( an 2 ann ) x2 ( an , n 1 ann ) xn 1
The system of equations in Eq. 5.68 can also be briefly described as follows:
r 1 b n aij ( r )
xi i j 1, j i x j , r 0,1, 2, n and i 1, 2, n (5.69)
aii aii
A sufficient condition for obtaining a solution by Jacobi’s iteration method is the diagonal
n
dominance, i.e. aii j 1, j i aij , i 1, 2, , n , in each row of A , modulus of the diagonal
element exceeds the sum of the off diagonal elements and also the diagonal elements, aii 0 . If
any diagonal element is 0, the equations can always be rearranged to satisfy this condition.
Gauss Seidel Iteration method: A simple modification to Jacobi’s iteration method is given by
Gauss – Seidel method. The steps are as follows:
Step 1: Determination of first approximation x1(1) , x2(1) , , xn(1) as,
(1) (0) (0) (0)
x1 (b1 a11 ) ( a12 a11 ) x2 ( a13 a11 ) x3 ( a1n a11 ) xn
x2(1) (b2 a22 ) ( a21 a22 ) x1(1) (a23 a22 ) x3(0) ( a2 n a22 ) xn(0)
(1) (1) (1) (0)
x3 (b2 a22 ) ( a21 a22 ) x1 ( a23 a22 ) x2 ( a2 n a22 ) xn (5.70)
xn(1) (bn ann ) ( an1 ann ) x1(1) ( an 2 ann ) x2(1) (an ,n 1 ann ) xn(1)1
Step ( + ): In general, if x1( n ) , x2( n ) , , xn( n ) are a system of nth approximations, then the next
approximation is given as,
( n 1) (n) ( n) ( n)
x1 (b1 a11 ) ( a12 a11 ) x2 (a13 a11 ) x3 ( a1n a11 ) xn
x2( n 1) (b2 a22 ) ( a21 a22 ) x1( n 1) ( a23 a22 ) x3( n ) ( a2 n a22 ) xn( n )
( n 1) ( n 1) ( n 1) (n)
x3 (b2 a22 ) ( a21 a22 ) x1 (a23 a22 ) x2 ( a2 n a22 ) xn (5.71)
xn( n 1) (bn ann ) ( an1 ann ) x1( n 1) ( an 2 ann ) x2( n 1) (an ,n 1 ann ) xn( n11)
The Eq. 5.71 can be briefly written as,
r 1 i 1 ( r 1) n (r )
xi (bi aii ) j 1, j i ( aij aii ) x j j i 1 ( aij aii ) x j , r 0,1, 2, n and i 1, 2, n (5.72)
Example 5.97 [MA-2014 (2 marks)]: Using the Gauss-Seidal iteration method with the initial guess
x (0)
1
3.5, x2(0) 2.25, x3(0) 1.625 , the second approximation x1(2) , x2(2) , x3(2) for the solution to the
system of equations: 2 x1 x2 7 , x1 2 x2 x3 1 , x2 2 x3 1 , is
(a) x1(2) 5.3125, x2(2) 4.4491, x3(2) 2.1563 (b) x1(2) 5.3125, x2(2) 4.3125, x3(2) 2.6563
(c) x1(2) 5.3125, x2(2) 4.4491, x3(2) 2.6563 (d) x1(2) 5.4991, x2(2) 4.4491, x3(2) 2.1563
2 x1 x2 0 x3 7
Solution (b): Let the system of equation be x1 2 x2 x3 1 and let the initial values be
0 x1 x2 2 x3 1
x1(0) 3.5, x2(0) 2.25, x3(0) 1.625 . So from Eq. 5.71, we have the first iteration as:
x1(1) ( x2(0) 7) 2 4.625 x1(2) ( x2(1) 7) 2 5.3125
(1) (1) (0) (2) (2) (1)
x2 (1 x1 x3 ) 2 3.625 ; so, 2nd iteration is x2 (1 x1 x3 ) 2 4.3125 .
x3(1) (1 x2(1) ) 2 2.3125 x3(2) (1 x2(2) ) 2 2.6563
Example 5.99 [EE-2011 (2 marks)]: Solution of the variables x1 and x2 for the following equations
is to be obtained by employing the Newton-Raphson iterative method. 10 x2 sin x1 0.8 0 ;
10 x22 10 x2 cos x1 0.6 0 . Assuming the initial valued x1 0.0 and x2 1.0 , the Jacobian matrix is
10 0.8 10 0 0 0.8 10 0
(a) (b) (c) (d)
0 0.6 0 10 10 0.6 10 10
Solution (b): As we have two equation and two unknowns so the Jacobian matrix is a 2 2 matrix
fi ( x) f1 (x) x1 f1 ( x) x2
whose elements are given by J , 1 i, j 2 . So J and
x j f 2 ( x) x1 f2 (x) x2
from the given data we have f1 ( x) 10 x2 sin x1 0.8 and f 2 ( x ) 10 x22 10 x2 cos x1 0.6 . So
10 x2 cos x1 10 sin x1 10 0
J at initial value x1 0.0 , x2 1.0 , we have J .
10 x2 sin x1 20 x2 10 cos x1 0 10
Example 5.100 [CS-1994 (2marks)]: Subjective Question: Match Column-I with Column-II
Column-I Column-II
A. Newton-Raphson 1. Integration
B. Runge-Kutta 2. Root finding
C. Gauss-Seidal 3. Ordinary Differential Equation
D. Simpson’s Rule 4. Solutions of Systems of Linear Equations
Solution: We know that (i) Newton-Raphson method is for Root finding; (ii) Runge-Kutta method is
for solving ODE; (iii) Gauss-Seidal method is for solving System of Linear Equations; (iv) Simpson’s
rule is for Integration. So answer is A – 2; B – 3; C – 4; D – 1.
Example 5.101 [EE-1997 (1 mark)]: Gauss-Seidel iterative method can be used for solving a set of
(a) linear differential equations only (b) both linear and non-linear algebraic equations
(c) linear algebraic equations only (d) both linear and non-linear differential equations
Solution (c): We know that Gauss-Seidel iterative method is used for solving a set of linear algebraic
equations only. So option (b) is correct.
Example 5.105 [MT-2010 (1 mark)]: Which of the following is an iterative technique to solve a
linear system of equations?
(a) Gaussian elimination (b) LU decomposition (c) Newton-Raphson (d) Jacobi method
Solution: We know that Gauss Elimination and LU decomposition method are used for solving
System of Linear Equations exactly; Newton-Raphson method is used for solving non-linear
equations; but Jacobi and Gauss-Seidel method are iterative methods which are used for solving linear
system of equations. So option (d) is correct.
Example 5.106 [MT-2011 (1 mark)]: Which one of the following methods is NOT used for
numerically solving an ordinary differential equation (ODE)?
(a) Euler’s method (b) Runge-Kutta method
(c) Adam-Bashforth method (d) Newton-Raphson method
Solution (d): As Euler’s, Runge-Kutta and Adam-Bashforth methods are used for solving an ODE;
but Newton-Raphson method is used for solving root finding. So option (d) is correct.
Example 5.107 [MT-2012 (1 mark)]: Which one of the following methods is NOT used for
integration?
(a) Rectangular rule (b) Trapezoidal rule (c) Simpson’s rule (d) Cramer’s rule
Solution (d): As Rectangular, Trapezoidal and Simpson’s rule are used for integration; but Cramer’s
rule is used for finding solution of system of linear equations having number of equations is equal to
number of unknowns. So option (d) is correct.
Example 5.108 [EC-2014 (1 mark)]: Match the application to appropriate numerical method.
Application Numerical Method
A. Numerical integration 1. Newton-Raphson method
B. Solution to transcendental equation 2. Runge-Kutta method
C. Solution to a system of linear equations 3. Simpson’s 1/3rd rule
D. Solution to a differential equation 4. Gauss elimination method
(a) A – 3, B – 2, C – 4, D – 1 (b) A – 3, B – 1, C – 4, D – 2
(c) A – 4, B – 1, C – 3, D – 2 (d) A – 2, B – 1, C – 3, D – 4
Solution (b): We know that (i) Simpson’s rule is for Numerical integration; (ii) Gauss Elimination
method is for solving System of Linear Equations; (iii) Runge-Kutta method is used for solving an
ODE; (iv) Newton-Raphson method is for solving non-linear algebraic or transcendental equations.
So option (b) is correct.
Example 5.109 [CH-2016 (1 mark)]: Which one of the following is an iterative technique for
solving a system of simultaneous linear algebraic equations?
(a) Gauss elimination (b) Gauss-Jordon (c) Gauss-Seidel (d) LU decomposition
Solution (c): Gauss Elimination, Gauss-Jordon and LU decomposition belong to the direct methods
for solving systems of linear equations; these are methods that yield solutions after an amount of
computations that can be specified in advance. But Gauss-Seidel is an iterative methods in which we
start from an initial value and obtain better and better approximations from a computational cycle
repeated as often as may be necessary. So option (c) is correct.
Exercise: 5.2
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
5
1. The value of 1 1 x 2 dx by using trapezoidal rule, with eight equal intervals, is _____.
2. The following points, given in the table, were found empirically. By using the trapezoidal rule the
3.6
value of 2.1 ydx is _____.
3. The difference between the exact value of 0 sin xdx and the approximate value of 0 sin xdx by
using trapezoidal rule, using 10 equal intervals, is
(a) 0.015 (b) 0.017 (c) 0.019 (d) 0.021
4. The vertical distance (in m) covered by a rocket from t 8 sec to t 30 sec is given by
30
8 2000 ln 140000 (140000 2100t ) 9.8t dt . Using single segment trapezoidal rule, the
distance covered (in m) is
(a) 11868 (b) 18168 (c) 16188 (d) 18618
5. The absolute relative true error for the distance covered in Q. 4 is _____.
2t , 1 t 5
10. The velocity of a body is given by v (t ) 2
, where t is given in seconds, v is
5t 3, 5 t 14
given in m/s. Using two-segment Simpson’s 1/3 rule, the distance in metres covered by the body
from t 2 to t 9 seconds most nearly is
(a) 949.33 (b) 1039.7 (c) 1200.5 (d) 1442.0
0.7
11. The value of 0.1 (1 x)dx , which 6 equal intervals is
(a) 1.86 (b) 1.97 (c) 1.78 (d) 1.90
12. The following data of the velocity of a body is given as a
function of time. The best estimate of the distance in metres Time 4 7 10 15
covered by the body from t 4 and t 15 using combined Velocity 22 24 37 46
Simpson’s 1/3 rule and the trapezoidal rule would be
(a) 354.70 (b) 362.50 (c) 368.00 (d) 378.80
13. The Velocity v (m/s) of a particle at distance from a
point on its path is given in the table. Estimate the time x: 0 10 20 30 40 50 60
(sec) taken to travel 60 m by using Simpson’s 3/8 rule. v : 47 58 64 65 61 52 38
(a) 1.064 (b) 1.054 (c) 1.404 (d) 1.304
14. Taking y (0) 0 and using forward Euler’s method with step size h 0.1 solution of the
differential equation dy dx 1 y gives the value of y (0.3) as _____.
15. Taking y (0) 1 and using modified Euler’s method with step size h 0.01 solution of the
differential equation dy dx x 2 y gives the value of y at x 0.02 as _____.
16. Using Runge-Kutta method, find an approximate value y at x 0.2 for the differential equation
dy dx x y , if y 1 when x 0 . _____.
17. The convergence of which of the following method is sensitive to starting value?
(a) False Position (b) Gauss Seidal method
(c) Newton-Raphson method (d) All (a), (b) and (c)
18. In the Gauss elimination method for solving a system of linear algebraic equations,
triangularization leads to _____ matrix.
(a) diagonal (b) lower triangular (c) upper triangular (d) singular
19. The rate of convergence of Gauss Seidal method is _____ that of Gauss Jacobi method.
(a) once (b) twice (c) thrice (d) reciprocal
20. _____ method is very fast compared to other methods.
(a) Gauss elimination (b) Gauss Jordon (c) Gauss Seidal (d) Gauss Jacobi
21. The lower triangular matrix [ L] in [ L][U ] decomposition of the matrix
25 5 4 1 0 0 u11 u12 u13
10 8 16 l 1 0 0 u22 u23 is
21
8 12 22 l31 l32 1 0 0 u33
1 0 0 1 0 0 1 0 0 1 0 0
(a) 0.40000 1 0 (b) 10 1 0
(c) 0.400 1 0 (d) 25 1 0
0.32000 1.7333 1 8 12 1 0.320 1.50 1 5 4 1
Answers Keys
Answer Keys: Exercise: 5.1
1 2 3 4 5 6 7 8 9 10 11 12 13 14 15
b b c a b a c a c d a b b a b
16 17 18 19 20 21 22 23 24 25 26 27 28 29 30
a d b b b c 2 c c 11 a a b c b
31 32 33 34 35 36 37 38 39
a d b c 36.85 0.5 b 4.75 c
[AG – 2017 (1 mark)]: A box contains three white and four red balls. Two balls are drawn randomly
in sequence. If the first draw resulted in a red ball, the probability of getting a second red ball in the
next draw is
(a) 0.33 (b) 0.50 (c) 0.67 (d) 0.75
Solution: The required situation is given in
following figure.
So required probability is given by:
P( R) P ( R R )
P( R R)
P ( R ) P ( R R ) P (W ) P ( R W )
(4 7) (3 6)
P( R R ) 0.5
(4 7) (3 6) (3 7) (4 6)
[AG – 2017 (1 mark)]: The probability of getting two heads and two tails from four tosses of the
same coin is _____.
Solution: The total number of cases when four coins are tossed is 24 16 .
Now we want 2H and 2T among 16 cases, which is equal to total number of different arrangement of
4! 6
H, H, T, T, which is 6 . Thus required probability is 0.375 .
2! 2! 16
[BT – 2017 (2 marks)]: The distribution of marks scored by a large class in an exam can be
represented as a normal distribution with mean and standard deviation . In a follow-up exam in
the same class, everyone scored 5 marks more than their respective score in the earlier exam. For this
follow-up exam, the distribution of marks can be represented as a normal distribution with mean 2
and standard deviation 2 . Which one of the following is correct?
(a) 2 ; 2 (b) 2 ; 2 (c) 2 ; 2 (d) 2 ; 2
x1 x2 x3 xn
Solution (d): As ; and
n
( x1 5) ( x2 5) ( x3 5) ( xn 5) x1 x2 x3 xn
2 5 5 . So 2 .
n n
Also n
i 1
( xi )
2
N ; and
2 n
i 1
{( xi 5) 2 }
2
N n
i 1
{( xi 5) ( 5)}
2
N n
i 1
{xi }
2
N .
[Similar question was also asked in MT – 2017 (1 mark)]
[BT – 2017 (2 marks)]: For the probability density P ( x) 0.5e 0.5 x , the integral 0 P( x ) dx _____.
Solution: 0 P ( x )dx 0.5e 0.5 x dx ( e 0.5 x ) 0 0 ( 1) 1 .
0
[CE – 2017 (1 mark)]: The number of parameters in the univariate exponential and Gaussian
distributions, respectively, are
(a) 2 and 2 (b) 1 and 2 (c) 2 and 1 (d) 1 and 1
e x , x 0, 0
Solution (b): For Exponential distribution, the distribution function is F ( x )
0, otherwise
in which unknown is . For Gaussian distribution, the distribution function is
2
x 1 1 x
F ( x) exp dx in which unknowns are and .
2
2
[CE – 2017 (1 mark)]: A two-faced coin has its faces designated as head (H) and tail (T). This coin is
tossed three times in succession to record the following outcomes: H, H, H. If the coin is tossed one
more time, the probability (up to one decimal place) of obtaining H again, given the previous
realizations of H, H and H, would be _____.
Solution: As first three outcomes are already Heads; now if the coin is tossed again then the outcome
does not depends on previous outcomes. So required probability of getting Head is 1 2 0.5 .
[CH – 2017 (1 mark)]: The marks obtained by a set of students are: 38, 84, 45, 70, 75, 60, 48. The
mean and median marks, respectively, are
(a) 45 and 75 (b) 55 and 48 (c) 60 and 60 (d) 60 and 70
38 84 45 70 75 60 48
Solution (c): Mean 60
7
Arranging the data in increasing order: 38, 45, 48, 60, 70, 75, 84; as we have n 7 , so median is
{(7 1) 2}th 4th observation, which is 60. So median is 60.
[CH – 2017 (2 marks)]: A box has 6 red balls and 4 white balls. A ball is picked at random and
replaced in the box, after which a second ball is picked. The probability of both the balls being red,
rounded to 2 decimal places, is _____.
6 6 9
Solution: The required probability is 0.36 .
10 10 25
[CS – 2017 (1 mark)]: Let X be a Gaussian random variable with mean 0 and variance 2 . Let
Y max( X , 0) where max( a, b) is the maximum of a and b . The median of Y is _____.
Solution: As ‘ X ’ is Gaussian random variable, so X N (0, 2 ) for X . Given
0, X 0
Y max( X , 0) . Since median is positional average, so median of Y is ‘0’.
X , 0 X
[CS – 2017 (2 marks)]: The number of integers between 1 and 500 (both inclusive) that are divisible
by 3 or 5 or 7 is _____.
Solution: Let n( I x ) be the number of all integers between 1 and 500 (both inclusive) that are
divisible by x . So n( I x ) I x , where denotes greatest integer function.
[CS – 2017 (2 marks)]: P and Q are considering to apply for a job. The probability that P applies
for the job is 1 4 , the probability that P applies for the job given that Q applies for the job is 1 2 ,
and the probability that Q applies for the job given that P applies for the job is 1 3 . Then the
probability that P does not apply for the job given that Q does not apply for the job is
(a) 4 5 (b) 5 6 (c) 7 8 (d) 11 12
Solution: Let A and B be the event which denote that P and Q applies for the job. So P ( A) 1 4 ;
P ( A B ) 1 2 ; P ( B A) 1 3 . As P ( A B ) P( B) P( A B ) P( A) P( B A) .
Thus P ( A B ) (1 4) (1 3) 1 12 P ( B ) (1 2) (1 4) (1 3) P ( B ) 1 6
P( A B ) 1 P( A B) 1 {P( A) P( B ) P ( A B)} 1 {(1 4) (1 6) (1 12)} 4
P( A B )
P( B ) 1 P( B ) 1 P( B ) 1 (1 6) 5
[CS – 2017 (2 marks)]: If a random variable X has a Poisson distribution with mean 5, then the
expression E[( X 2) 2 ] equals _____.
Solution: E[( X 2) 2 ] E ( X 2 4 X 4) E ( X 2 ) 4 E ( X ) 4
For Poisson distribution, E ( X 2 ) a 2 a ; and E ( X ) a mean of Poisson distribution 5 .
So E[( X 2) 2 ] a 2 a 4a 4 a 2 5a 4 25 25 4 54 .
[EC – 2017 (1 mark)]: Three fair cubical dice are thrown simultaneously. The probability that all
three dice have the same number of dots on the faces showing up is (up to third decimal place) _____.
Solution: Total number of possible outcomes when three fair cubical dice are thrown simultaneously
is 63 216 . Total number of favourable cases: {(1,1,1), (2, 2, 2), , (6, 6, 6)} 6 cases. So required
probability is 6 63 0.0277 .
[EC – 2017 (2 mark)]: Passengers try repeatedly to get a seat reservation in any train running
between two stations until they are successful. If there is 40% chance of getting reservation in any
attempt by a passenger, then the average number of attempts that passengers need to make to get a
seat reserved is _____.
Solution: Let X be a random variable which takes number of attempts. Given probability of any
attempt to be successful is p 40% 40 100 2 5 0.4 ; so probability of any attempt to be
unsuccessful is q 1 p 0.6 .
X 1 2 3 4
P( X ) p q p qq p qqq p
E ( X ) XP( X ) 1 p 2( q p ) 3( q q p ) 4( q q q p ) p(1 2q 3q 2 4q 3 )
E ( X ) p (1 2q 3q 2 4q 3 ) …(i)
On multiplying (i) with q , we get qE ( X ) p ( q 2q 2 3q 3 4q 4 ) …(ii). So (i) – (ii)
p p 0.4 1
(1 q) E ( X ) p (1 q q 2 q 3 ) E( X ) 2
2
2.5
1 q (1 q) (1 0.6) 0.4
[EE – 2017 (1 mark)]: An urn contains 5 red balls and 5 black balls. In the first draw, one ball is
picked at random and discarded without noticing its colour. The probability to get a red ball in the
second draw is
(a) 1 2 (b) 4 9 (c) 5 9 (d) 6 9
Solution: Let ‘ R ’ be red ball; and ‘ B ’ be black ball. The
required situation is given in figure. So the total
probability for getting a red ball in the second draw is
P ( R ) P ( R R) P ( B ) P ( R B)
5 4 5 5 45
0.5
10 9 10 9 90
[EE – 2017 (2 marks)]: A person decides to toss a fair coin repeatedly until he gets a head. He will
make at most 3 tosses. Let the random variable Y denote the number of heads. The value of var{Y } ,
where var{} denotes the variance, equals
(a) 7 8 (b) 49 64 (c) 7 64 (d) 105 64
Solution (c): Let X be a random variable which takes X 1 1 1
number of attempts. The probability of getting head is P( X ) p q p qq p
p 1 2 ; and probability of getting tail is q 1 2 . Now
E ( X ) XP( X ) 1 p 1 ( q p ) 1 ( q q p ) p (1 q q 2 ) (1 2){1 (1 2) (1 2) 2 }
E( X ) 7 8
E ( X 2 ) X 2 P ( X ) 12 p 12 (q p ) 12 (q q p ) p (1 q q 2 )
E ( X 2 ) (1 2){1 (1 2) (1 2) 2 } 7 8
So Var ( X ) E ( X 2 ) {E ( X )}2 (7 8) (7 8) 2 7 64 .
[MA – 2017 (2 marks)]: Let E and F be any two events with P ( E ) 0.4 , P ( F ) 0.3 and
P ( F | E ) 3P( F | E c ) . Then P ( E | F ) equals (rounded to 2 decimal places) _____.
P( F E ) P( F E c ) P( F E ) P( F ) P( F E )
Solution: P ( F | E ) 3P ( F | E c ) 3 c
3
P( E ) P( E ) P( E ) 1 P( E )
P( F E ) 0.3 P ( F E ) P( F E ) 0.2
3 P ( F E ) 0.2 . So P ( E | F ) 0.66 .
0.4 1 0.4 P( F ) 0.3
[ME – 2017 (1 mark)]: A six-face fair dice is rolled a large number of times. The mean value of the
outcomes is _____.
Solution: Let X be a random variable which takes number appear on dice; and P ( X ) be the
probability of appearance of that number. So
X 1 2 3 4 5 6
P( X ) 16 16 16 16 16 16
Mean E ( X ) XP( X ) 1(1 6) 2(1 6) 3(1 6) 4(1 6) 5(1 6) 6(1 6) 3.5
[IN – 2017 (2 marks)]: The probability that a communication system will have high fidelity is 0.81.
The probability that the system will have both high fidelity and high selectivity is 0.18. The
probability that a given system with high fidelity will have high selectivity is
(a) 0.181 (b) 0.191 (c) 0.222 (d) 0.826
Solution (c): Let A and B be the events respectively that the communication system will have high
fidelity and high selectivity, respectively, then P ( A) 0.81 and P ( A B) 0.18 . So required
probability that a given system with high fidelity will have high selectivity is given by
P ( B A) P ( A B ) P ( A) 0.18 0.81 0.222 .
[ME – 2017 (1 mark)]: Two coins are tossed simultaneously. The probability (up to two decimal
point accuracy) of getting at least one head is _____
2
Solution: Total number of possible outcomes is 2 4 . Set of favourable outcomes is
{( H , T ), (T , H ), ( H , H )} ; so total number of favourable outcomes is 3 . So required probability is
3 4 0.75 .
[ME – 2017 (1 mark)]: A sample of 15 data is as follows: 17, 18, 17, 17, 13, 18, 5, 5, 6, 7, 8, 9, 20,
17, 3. The mode of the data is
(a) 4 (b) 13 (c) 17 (d) 20
Solution (c): Mode is that value of data whose frequency of occurrence is most. So mode for the
given data is 17.
[MN – 2017 (1 mark)]: F ( y ) and f ( y ) are the probability distribution function and density
function respectively of a continuous variable Y in the interval (0, ) . Which one of the following is
TRUE?
y
(a) F ( y ) f ( x ) dx (b) F ( y ) f ( x) dx (c) F ( y ) {df ( y)} dy (d) F ( y ) 1 f ( y )
y 0
Solution (b): We know that if X be a continuous random variable defined on a sample space S with
a
probability density function f ( x ) . Then cumulative distribution function F ( a ) f ( x ) dx .
In given problem we have a y and y (0, ) ; so
y 0 y y y
F ( y) f ( x )dx f ( x ) dx f ( x ) dx 0 f ( x ) dx f ( x) dx
0 0 0
[MN – 2017 (2 marks)]: In a book of 600 pages, there are 60 typographical errors. Assuming Poisson
distribution for the number of errors per page, the probability of no errors in randomly chosen 4 pages
is _____.
Solution: In the book per page there are 60 600 0.1 errors. So for 4 pages, there will be
4 0.1 0.4 errors, which is the mean of given Poisson distribution. Now
0 0.4
P ( x 0) 0.4 e 0! 0.67 .
[MT – 2017 (2 marks)]: Assume that the probability of South Africa winning against India is 1 3 . If
South Africa plays a 3 match cricket series against India, the probability that South Africa wins only
one match is (assume that no matches are tied or abandoned) _____ (answer up to three decimal
places)
Solution: Among three matches, South Africa can win only one match in one of the possibility:
1 2 2 2 1 2 2 2 1 4
{WLL, LWL, LLW } . So required probability is 0.44 .
3 3 3 3 3 3 3 3 3 9
ax 2 , 0 x5
[PE – 2017 (2 marks)]: The value of constant a for which: f ( x ) is a valid
0, otherwise
probability density function, is (given, a 0 ):
(a) 1 125 (b) 3 125 (c) 6 125 (d) 9 125
Solution (b): For continuous probability density function f ( x ) , we must have
0 5 5
f ( x)dx 1 f ( x)dx 0 f ( x) dx f ( x )dx 1 0 ax 2 0 1 a 3 125 .
5 0
[PI – 2017 (2 marks)]: Two machines are defective in a lot of 10. A combination of four machines is
to be picked at a time from the lot. The maximum number of combinations that can be obtained
without any defective machine is _____.
Solution: Total number of non-defective machines is 10 2 8 . From these 8 machines we have to
8!
choose 4 machines, which can be done in 8 C4 70 ways.
4! 4!
[TF – 2017 (2 marks)]: The median, accurate to two decimal places, of a random variable X with
3 x 4 , x 1
the probability density function f ( x ) is _____.
0, x 1
Solution: We know that a median of a quantity x distributed through a population is a value T such
that half of the population has values of x less than T and half the population has values of x greater
T 1
than T . So T satisfies f ( x ) .
2
3 x 4 , x 1
Now we have f ( x) , so if T is the median, then
0, x 1
T 1 1 T 4 1
f ( x) 2 0 dx 1 3x dx T 3 2 T 1.2599 1.26 .
2
[XE – 2017 (1 mark)]: Three fair dice are rolled simultaneously. The probability of getting a sum of
5 is
(a) 1 108 (b) 1 72 (c) 1 54 (d) 1 36
Solution: If three fair dice are rolled then total number of possible cases is 63 216 . Let x, y , z be
the number appearing on three dices, respectively. So total number of favourable cases will be total
number of triplets ( x, y , z ) such that x y z 5 , which is 51 C31 4 C2 6 . So required
probability is 6 63 1 36 .
6.1.1 Permutation
The Factorial: Let n be a ve integer, then continued product of first n natural numbers is
called factorial n , denoted by n! or n . Also, 0! 1 . Thus
n ! n( n 1)(n 2) 3 2 1 , n ! n (n 1)! n( n 1) (n 2)! and so on (6.1)
If n is ve integer or a fraction, n! is not defined.
Multiplication principle: Suppose that an event X can be decomposed into two stages A and
B . Let stage A occur in m ways and suppose that these stages are unrelated, in the sense that
stage B occurs in n ways regardless of the outcome of stage A . Then event X occur in mn
ways. This rule is applicable even if event X can be decomposed in more than two stages.
In other words, if one operation can be performed independently in m different ways and if
second operation can be performed independently in n different ways and a third operation
can be performed independently in p different ways and so on, then the total number of ways
in which all the operations can be performed in the stated order is ( m n p )
Example 6.1: A college offers 7 courses in the morning and 5 in the evening. Find the number of
ways a student can select exactly one course, either in the morning or in the evening.
Solution: The student has seven choices from the morning courses out of which he can select one
course in 7 ways. For the evening course, he has 5 choices out of which he can select one course in 5
ways. Hence he has total number of 7 5 12 choices.
Example 6.2 [CS-2004 (1 mark)]: The number of different n n symmetric matrices with each
x
element being either 0 or 1 is: (Note power (2, x) is same as 2 )
(a) power (2, n) (b) power (2, n 2 ) (c) power {2, ( n 2 n) 2} (d) power {2, ( n 2 n) 2}
Solution: The total number of elements in n n matrix is n 2 . In this matrix we have n number of
elements along the diagonal and ( n2 n) 2 number of elements above and below diagonal of the
matrix. If this matrix to be symmetric with each element being either 0 or 1 then we fill the diagonal
places and one of the places in above or below diagonal (note that if above diagonal elements are
filled then below diagonal elements will be adjusted in order to it have symmetric matrix). So we have
to fill n ( n 2 n) 2 ( n 2 n) 2 number of places with either 0 or 1. As each place have two choices
2 ( n2 n) 2
0 or 1; so ( n2 n) 2 number of places will be filled in 2 2 {( n n) 2 times} 2 ways
which is the number of different n n symmetric matrices with each element being either 0 or 1.
Example 6.4 [XE-2007 (2 marks)]: The number of n n matrices that are simultaneously Hermitian,
unitary and diagonal is
(a) 2
n
(b) n 2 (c) 2n (d) 2
Solution (a): Let U A iB be a matrix. For U to be diagonal matrix, both A and B must be
diagonal, i.e., A diag ( a1 , a2 , , an ) and B diag (b1 , b2 , , bn ) . For U to be hermitian matrix,
a j ib j a j ib j b j 0 B O and thus U A diag ( a1 , a2 , , an ) is real matrix. For U to
be unitary matrix (i.e., orthogonal), each row must be of unit length, i.e., we must have
a 2j 1 a j 1 . Now we have n places and each place has 2 options ( 1 or 1) to be filled; so
total number of ways is 2 2 2 (n times) 2 n .
Example 6.5 [CS-2014 (2 marks)]: The number of distinct positive integral factors of 2014 is …
Solution: As 2014 1 2 19 53 ; so positive integral factors of 2014 are 1, 2, 19, 53, 2 19 , 2 53
, 19 53 , 2014 . So number of distinct positive integral factors of 2014 is 8.
The number of arrangements that can be formed using n objects out of which p are identical
(and of one kind) q are identical (and of another kind), r are identical (and of another kind)
and the rest are distinct is ( n !) {( p !)( q !)(r !)} .
Example 6.6: How many numbers can be made with the help of the digits 0, 1, 2, 3, 4, 5 which are
greater than 3000 (repetition is not allowed)?
Solution: All the 5 digit numbers and 6 digit numbers are greater than 3000. Therefore number of 5
digit numbers is 6 P5 5 P5 600 (Since the case that 0 will be at ten thousand place should be omit).
Similarly number of 6 digit numbers 6! 5! 600 . Now the numbers of 4 digit numbers which are
greater than 3000, having 3, 4 or 5 at first place, this can be done in 3 ways and remaining 3 digit may
be filled from remaining 5 digits i.e., required number of 4 digit numbers are 5 P3 3 180 . Hence
total required number of numbers 600 600 180 1380 .
Example 6.7 [CS-2003 (1 mark)]: n couples are invited to a party with the condition that every
husband should be accompanied by his wife. However, a wife need not be accompanied by her
husband. The number of different gathering possible at the party is
(a) 2 n Cn 2 n (b) 3n (c) (2n)! 2n (d) 2n Cn
Solution (b): As for every couple, they have three options (i) nobody goes to gathering; (ii) wife
alone goes; (iii) both husband and wife goes. As there are n couples, so total number of ways of
gathering is 3 3 ( n times) 3n .
Example 6.8 [CH-2012 (1 mark)]: A box containing 10 identical compartments has 6 red balls and 2
blue balls. If each compartment can hold only one ball, then the number of different possible
arrangements are
(a) 1026 (b) 1062 (c) 1260 (d) 1620
Solution: We have total 8 balls (out of which 6 are of one type and 2 are of second type) and 10
places. So total number of different possible arrangements are (10!) {(6!)(2!)} 1260 .
Example 6.9 [MN-2013 (1 mark)]: The number of ways in which the letters in the word MINING
can be arranged is
(a) 90 (b) 180 (c) 360 (d) 720
Solution: As there are 6 letters in MINING; but two letters ‘I’ and ‘N’ appear two times so the
number of ways in which the letters in the word MINING can be arranged is (6!) {(2!)(2!)} 180 .
Example 6.10: In a train 5 seats are vacant then how many ways can three passengers sit?
Solution: Number of ways are 5 P3 (5!) {(5 3)!} (5!) (2!) 120 2 60 .
Example 6.11: The number of 5 digit telephone numbers having at least one of their digits repeated is
Solution: Using the digits 0,1, 2, , 9 the number of five digit telephone numbers which can be
formed is 105 (since repetition is allowed). The number of five digit telephone numbers which have
none of the digits repeated is 10 P5 30240 . Hence, the required number of telephone numbers is
5
10 30240 69760 .
Conditional Permutations
Number of permutations of n dissimilar things taken r at a time when p particular things
n p
always occur is Cr p r ! .
Number of permutations of n dissimilar things taken r at a time when p particular things never
n p
occur is Cr r !
The total number of permutations of n different things taken not more than r at a time, when
each thing may be repeated any number of times, is {n( n r 1)} ( n 1) .
Number of permutations of n different things, taken all at a time, when m specified things
always come together is ( m !) ( n m 1)! .
Number of permutations of n different things, taken all at a time, when m specified things never
come together is ( n !) (m !) ( n m 1)! .
Let there be n objects, of which m objects are alike of one kind, and the remaining (n m)
objects are alike of another kind. Then, the total number of mutually distinguishable permutations
that can be formed from these objects is ( n !) {( m !) ( n m)!} . This theorem can be extended
further i.e., if there are n objects, of which p1 are alike of one kind; p2 are alike of another kind;
th
p3 are alike of 3rd kind; … ; pr are alike of r kind such that p1 p2 pr n ; then the
number of permutations of these n objects is ( n !) {( p1 !) ( p2 !) ( pr !)} .
Gap method: Suppose 5 males A, B, C, D, E are arranged in a row as A B C D E .
There will be six gaps between these five. Four in between and two at either end. Now if three
females P, Q, R are to be arranged so that no two are together we shall use gap method i.e.,
arrange them in between these 6 gaps. Hence the answer will be 6 P3 .
Together: Suppose we have to arrange 5 persons in a row which can be done in 5! 120
ways. But if two particular persons are to be together always, then we tie these two particular
persons with a string. Thus we have 5 2 1 (1 corresponding to these two together)
3 1 4 units, which can be arranged in 4! ways. Now we loosen the string and these two
particular can be arranged in 2! ways. Thus total arrangements 24 2 48 . Never together
Total – Together 120 48 72 .
Example 6.12: Find the number of words which can be made out of the letters of the word
‘MOBILE’ when consonants always occupy odd places.
Solution: The word ‘MOBILE’ has three even places and three odd places. It has 3 consonants and 3
vowels. In three odd places we have to fix up 3 consonants which can be done in 3 P3 ways. Now
3
remaining three places we have to fix up remaining three places which can be done in P3 ways.
Hence, the total number of ways 3 P3 3 P3 36 .
Example 6.13: An n digit number is a positive number with exactly n digits. Nine hundred distinct
n digit numbers are to be formed using only the three digits 2, 5 and 7. Find the smallest value of n
for which this is possible.
Solution: Since at any place, any of the digits 2, 5 and 7 can be used total number of such positive
n digit numbers are 3n . Since we have to form 900 distinct numbers, hence 3n 900 n 7 .
Circular Permutations: So far we have been considering the arrangements of objects in a line.
Such permutations are known as linear permutations. Instead of arranging the objects in a line, if we
arrange them in the form of a circle, we call them, circular permutations. In circular permutations,
what really matters is the position of an object relative to the others. Thus, in circular permutations,
we fix the position of the one of the objects and then arrange the other
objects in all possible ways. There are two types of circular
permutations:
1. The circular permutations in which clockwise and the anticlockwise
arrangements give rise to different permutations, e.g. Seating
arrangements of persons round a table.
2. The circular permutations in which clockwise and the anticlockwise Figure 6.1: Circular
arrangements give rise to same permutations, e.g. arranging some Permutation
beads to form a necklace.
Suppose A, B, C , D are the four beads forming a necklace. As shown in Fig. 6.1, they have been
arranged in clockwise and anticlockwise directions in the first and second arrangements respectively.
Now, if the necklace in the first arrangement be given a turn, from clockwise to anticlockwise, we
obtain the second arrangement. Thus, there is no difference between the above two arrangements.
Example 6.14: In how many ways can 5 boys and 5 girls sit in a circle so that no boys sit together?
Solution: Since total number of ways in which boys can occupy any place is (5 1)! 4! and the 5
girls can be sit accordingly in 5! ways. Hence required number of ways are (4!) (5!) .
Example 6.15: Find the number of ways in which 5 beads of different colours form a necklace?
Solution: The number of ways in which 5 beads of different colours can be arranged in a circle to
form a necklace are (5 1)! 4! . But the clockwise and anticlockwise arrangement are not different
(because when the necklace is turned over one gives rise to another). Hence the total number of ways
of arranging the beads is (1 2)(4!) 12 .
6.1.2 Combination
Each of the different groups or selections which can be formed by taking some or all of a number of
objects, irrespective of their arrangements, is called a combination. Suppose we want to select two out
of three persons A , B and C . We may choose AB or BC or AC . Clearly, AB and BA represent
the same selection or group but they give rise to different arrangements. Clearly, in a group or
selection, the order in which the objects are arranged is immaterial.
n
Number of all combinations of n things, taken r at a time is denoted by C ( n, r ) or n C r or .
r
Difference between a permutation and combination
In a combination only selection is made whereas in a permutation not only a selection is made
but also an arrangement in a definite order is considered.
In a combination, the ordering of the selected objects is immaterial whereas in a permutation,
the ordering is essential. For example A , B and B , A are same as combination but different
as permutations.
Practically to find the permutation of n different items, taken r at a time, we first select r
items from n items and then arrange them. So usually the number of permutations exceeds
the number of combinations.
Each combination corresponds to many permutations. For e.g., the six permutations ABC ,
ACB , BCA , BAC , CBA and CAB correspond to the same combination ABC .
Generally we use the word ‘arrangements’ for permutations and word “selection” for
combinations.
Example 6.16 [CS-2004 (2 marks)]: Mala has a colouring book in which each English letter is drawn
two times. She wants to paint each of these 52 prints with one of k colours, such that the colour pairs
used to colour any two letters are different. Both prints of a letter can also be coloured with the same
colour. What is the minimum value of k that satisfies this requirement?
(a) 9 (b) 8 (c) 7 (d) 6
Solution (c): We have letters A-Z and each letter is printed twice, so there are 52 letters. Now we
have to colour each letter, so we need a pair of colours for that, because each letter is printed twice.
Also in a pair, both colours can be some. Now condition is that a pair of colours cannot be used more
than once. So let she has 2 colours: Black and White. She can colour as follows : (A,A): (Black,
Black), (B,B): (White, White), (C,C): (Black, White). Now we do not have more pairs of colours left,
we have used all pairs, but could colour only 3 letters out of 26. So our task is to find minimum no. of
colours, so that we could colour all 26 letters. So if she has k colours, she can have k pairs of same
colours, thus colouring k letters, then k C2 other pairs of colours, thus colouring k C2 more letters. So
total no. of letters colored = k k C2 k (1 2) k ( k 1) (1 2)k ( k 1) . So we want
2
(1 2) k (k 1) 26 k k 52 k ( , 7.7) [6.7, ) . As k cannot be negative so
k [6.7, ) . Thus minimum value of k that satisfies the requirement is 7.
Example 6.17 [CS-2004 (2 marks)]: In an m n matrix such that all non-zero entries are covered in
a rows and b columns. Then the maximum number of non-zero entries, such that no two are on the
same row or column, is
(a) a b (b) max( a, b) (c) min( m a, n b) (d) min( a, b)
Solution (d): Suppose a b , for example let a 3 , b 4 , then we can put non-zero entries only in 3
rows and 4 columns. So suppose we put non-zero entries in any 3 rows in 3 different columns. Now
we cannot put any other non-zero entry anywhere in matrix, because if we put it in some other row,
then we will have 4 rows containing non-zeros, if we put it in one of those 3 rows, then we will have
more than one non-zero entry in one row, which is not allowed. So we can fill only ‘ a ’ non-zero
entries if a b , similarly if b a , we can put only ‘ b ’ non-zero entries. So answer is min( a, b) ,
because whatever is less between a and b , we can put at most that many non-zero entries.
Statement for Linked Answer Questions 6.18 & 6.19: Suppose that a robot is placed on the
Cartesian plane. At each step it is allowed to move either one unit up or one unit right, i.e., if it is at
(i, j ) then it can move to either (i 1, j ) or (i, j 1) .
Example 6.18 [CS-2007 (2 marks)]: How many distinct paths are there for the robot to reach the
point (10,10) starting from the initial position (0, 0)
(a) 20 C10 (b) 2
20
(c) 2
10
(d) None of these
Solution (a): At each move, robot can move either 1 unit right, or 1 unit up, and there will be 20 such
moves required to move from (0,0) to (10,10). So we have to divide these 20 moves, numbered from 1
to 20, into 2 groups: right group and up group. Each group contains 10 elements each. So basically,
we have to divide 20 things into 2 groups of 10 and 10 things each, which can be done in
(20!) {(10!)(10!)} 20 C10 .
Example 6.19 [CS-2007 (2 marks)]: Suppose that the robot is not allowed to traverse the line
segment from (4, 4) to (5, 4) . With this constraint, how many distinct paths are there for the robot to
reach (10,10) starting from (0, 0) ?
(a) 2
9
(b) 2
19 (c) 8 C4 11C5 (d) 20 C10 8C4 11C5
Solution (d): Since we are not allowed to traverse from (4,4) to (5,4). So we will subtract all those
paths which were passing through (4,4) to (5,4). To count number of paths passing through (4,4) to
(5,4), we find number of paths from (0,0) to (4,4), and then from (5,4) to (10,10). From (0,0) to (4,4),
number of paths is (8!) {(4!)(4!)} 8C4 (found in same way as earlier). From (5,4) to (10,10),
number of paths is (11!) {(5!)(6!)} 11C5 . So total number of paths required is 20
C2 8C4 11C5 .
Example 6.20 [CS-2006 (2 marks)]: Given a set of elements N {1, 2, , n} and two arbitrary
subsets A N and B N , how many of the n! permutations from N to N satisfy
min ( A) min ( B ) , where min( S ) is the smallest integer in the set of integers S , and ( S ) is
the set of integers obtained by applying permutation to each elements of S ?
A B n
(a) n A B A B
2
(b) A B n 2 2
(c) n !
A B
(d) A B
2
A B
Solution (c): First let us understand what question is asking. So is a function from N to N ,
which just permutes the elements of N , so there will be n! such permutations. Now given a
particular , i.e., given a particular permutation scheme, we have to find number of permutations out
of these n! permutations in which minimum elements of A and B after applying to them are
same. So for example, if N {1, 2, 3} , is {2, 3,5} , and if A is {1, 3} , then ( A) {1, 2} . Now
n
number of elements in A B is A B . We can choose permutations for A B in C A B ways.
Note that here we are just choosing elements for permutation, and not actually permuting. Let this
chosen set be P . Now once we have chosen numbers for permutations, we have to select mapping
from each element of A B to some element of P . So first of all, to achieve required condition
specified in question, we have to map minimum number in P to any of the number in A B , so that
min{ ( A)} min{ ( B )} . We can do this in A B ways, as we can choose any element of A B
to be mapped to minimum number in P . Now we can permute numbers in P in A B 1 ! ways,
since one number (minimum) is already fixed. Also, we can permute remaining in n A B 1 !
n
ways, so total no. of ways is C A B A B A B 1 ! n A B 1 ! n ! A B A B .
Example 6.22: How many words of 4 consonants and 3 vowels formed from 6 consonants and 5
vowels.
Solution: Required number of words = 6 C4 5C3 7! = 756000
Example 6.23: In a conference of 8 persons, if each person shake hand with the other one only, then
find the total number of shake hands.
Solution: Total number of shake hands when each person shake hands with the other once only is
8
C2 28 ways.
Example 6.24 [CS-1999 (2 marks)]: Two girls have picked 10 roses, 15 sunflowers and 14 daffodils.
What is the number of ways they can divide the flowers amongst themselves?
(a) 1638 (b) 2100 (c) 2640 (d) None of these
Solution (c): Two girls can pick roses in 11 ways with in which one girl gets 0,1, 2, ,10 number of
roses and the other girl gets the remaining. Similarly two girls can pick sunflower in 16 ways with in
which one girl gets 0,1, 2, ,15 number of sunflowers and the other girl gets the remaining. Similar
case happens for the daffodils, i.e., for daffodil flowers 15 ways
So total number of ways they can divide the flowers amongst themselves is 11 16 15 2640 ways.
Example 6.25: A man has 10 friends. In how many ways he can invite one or more of them to a
party?
10
Solution: Required number of friend 2 1 (As the case that no friend be invited is excluded)
Example 6.26: Find the numbers greater than 1000 but not greater than 4000 which can be formed
with the digits 0, 1, 2, 3, 4 (repetition of digits is allowed).
Solution: Numbers greater than 1000 and less than or equal to 4000 will be of 4 digits and will have
either 1 (except 1000) or 2 or 3 in the first place with 0 in each of remaining places. After fixing 1st
place, the second place can be filled by any of the 5 numbers. Similarly third place can be filled up in
5 ways and 4th place can be filled up in 5 ways. Thus there will be 5 5 5 125 ways in which 1 will
be in first place but this include 1000 also hence there will be 124 numbers having 1 in the first place.
Similarly 125 for each 2 or 3. One number will be in which 4 in the first place and i.e., 4000. Hence
the required numbers are 124 125 125 1 375 ways.
Conditional Combination
The number of ways in which r objects can be selected from n different objects if k particular
objects are
Always included n k C r k Never included n k C r
n p
C r n p C r 1 n p C r 2 n p C0 if r p ; and
n p
Cr n p Cr 1 n p Cr 2 n p Cr p if r p
Example 6.27 : In the 13 cricket players 4 are bowlers, then how many ways can form a cricket team
of 11 players in which at least 2 bowlers included?
Solution: The number of ways can be given as follows: 2 bowlers and 9 other players 4 C2 9C9 ; 3
bowlers and 8 other players 4C3 9C8 ; 4 bowlers and 7 other players 4 C4 9C7 . Hence required
number of ways 6 1 4 9 1 36 78 .
Example 6.28: In how many ways a team of 10 players out of 22 players can be made if 6 particular
players are always to be included and 4 particular players are always excluded?
Solution: 6 particular players are always to be included and 4 are always excluded, so total number of
selection, now 4 players out of 12. Hence number of ways 12 C4 .
Example 6.29: In how many ways can 5 prizes be distributed among four students when every
student can take one or more prizes?
Solution: The required number of ways 45 1024 [since each prize can be distributed by 4 ways]
Example 6.30: Find the number of ways in which 9 persons can be divided into three equal groups?
Solution: Total ways {9!} {(3!)3 } 280 .
De-arrangement: Any change in the given order of the things is called a derangement. If n things
form an arrangement in a row, the number of ways in which they can be deranged so that no one of
1 1 1 1
them occupies its original place is n ! 1 ( 1) n .
1! 2! 3! n!
Example 6.31: There are four balls of different colours and four boxes of colours same as those of the
balls. Find the number of ways in which the balls, one in each box, could be placed such that a ball
doesn't go to box of its own colour.
1 1 1 1
Solution: Number of derangement are 4! 1 12 4 1 9 .
1! 2! 3! 4!
Example 6.33: Out of 18 points in a plane, no three are in the same straight line except five points
which are collinear. Find the number of (i) straight lines (ii) triangles which can be formed by joining
them.
Solution: Out of 18 points, 5 are collinear.
(i) Number of straight lines 18C2 5C2 1 153 10 1 144
(ii) Number of triangles 18 C3 5C3 816 10 806
Multinomial Theorem: Let x1 , x2 , , xm be integers. Then (i) number of solutions to the equation
x1 x2 xm n (ii) subject to the condition a1 x1 b1 , a2 x2 b2 , , am xm bm is equal
a a 1 b a a 1 b a
to (iii) the coefficient of x n in ( x 1 x 1 x 1 ) ( x 2 x 2 x 2 ) ( x am x m1 xbm ) .
This is because the number of ways, in which sum of m integers in (i) equals n , is the same as the
number of times x n comes in (iii).
Use of solution of linear equation and coefficient of a power in expansions to find the number of
ways of distribution:
The number of integral solutions of x1 x2 x3 xr n where x1 0, x2 0, xr 0 is the
same as the number of ways to distribute n identical things among r persons. This is also equal
to the coefficient of x n in the expansion of ( x 0 x1 x 2 x 3 ) r
( r n 1)!
coefficient of x n in {1 (1 x )}r coefficient of x n in (1 x ) r n r 1Cr 1
n !( r 1)!
The number of integral solutions of x1 x2 x3 xr n where x1 1, x2 1, , xr 1 is
same as the number of ways to distribute n identical things among r persons each getting at least
1. This also equal to the coefficient of x n in the expansion of ( x1 x 2 x 3 ) r coefficient of
( n 1)!
n r n r
x in { x (1 x )} coefficient of x in x (1 x)
r
n 1Cr 1
( n r )!( r 1)!
Example 6.34 [CS-2003 (2 marks)]: m identical balls are to be placed in n distinct bags. You are
given that m kn , where k is a natural number 1 . In how many ways can the balls be placed in the
bags if each bag must contain at least k balls?
mk m kn n 1 m 1 m kn n k 2
(a) (b) (c) (d)
n 1 n 1 nk nk
Solution (b): Since we want at least k balls in each bag, so first we put kn balls into bags, i.e., k
balls in each bag. Now we are left with m kn identical balls, and we have to put them into n distinct
bags such that each bag may receive zero or more balls. So applying the above theorem, we get
number of ways to be m kn n 1 Cn 1 .
Example 6.35: A student is allowed to select utmost n books from a collection of (2n 1) books. If
the total number of ways in which he can select one book is 63, then find the value of n .
Solution: As the student is allowed to select utmost n books out of (2n 1) books; in order to
select one book he has the choice to select one, two, … , n books. If T is the total number of ways of
selecting one book T 2 n 1C1 2 n 1C 2 2 n 1Cn 63 . So, the sum of binomial coefficients,
2 n 1C0 2 n 1C1 2 n 1C2 2 n 1Cn 2 n 1Cn 1 2 n 1Cn 2 2 n 1C2 n 1 (1 1) 2 n 1 22 n 1
2 n 1C0 2( 2 n 1 C1 2 n 1C2 2 n 1Cn ) 2 n 1C2 n 1 22 n 1
1 2(T ) 1 22 n 1 1 T (1 2)22 n 1 22 n 1 63 2 2 n 2 6 2 2 n n 3 .
Number of Divisors: Let N p11 p 2 2 p3 3 pk k , where p1 , p2 , , pk are different primes and
1 , 2 , , k are natural numbers then:
The total number of divisors of N including 1 and N is (1 1) ( 2 1) ( k 1)
The total number of divisors of N excluding 1 and N is (1 1) ( 2 1) ( k 1) 2
The total number of divisors of N excluding 1 or N is (1 1) ( 2 1) ( k 1) 1
The sum of these divisors is ( p10 p12 p11 )( p20 p21 p22 ) ( pk0 p1k pk k )
Important Points
All the numbers whose last digit is an even number 0, 2, 4, 6 or 8 are divisible by 2.
All the numbers sum of whose digits are divisible by 3, is divisible by 3, e.g. 534. Sum of the
digits is 12, which are divisible by 3, and hence 534 is also divisible by 3.
All those numbers whose last two-digit number is divisible by 4 are divisible by 4, e.g. 7312,
8936, are such that 12, 36 are divisible by 4 and hence the given numbers are also divisible by 4.
All those numbers, which have either 0 or 5 as the last digit, are divisible by 5.
All those numbers, which are divisible by 2 and 3 simultaneously, are divisible by 6, e.g., 108.
All those numbers whose last three-digit number is divisible by 8 are divisible by 8.
All those numbers sum of whose digit is divisible by 9 are divisible by 9.
All those numbers whose last two digits are divisible by 25 are divisible by 25, e.g., 73125.
Example 6.36: Find the total number of divisors of 9600 including 1 and 9600.
Solution: Since 9600 2 7 31 52 . Hence number of divisors (7 1)(1 1)(2 1) 48 .
Distributing balls into different boxes: In this section, we want to consider the problem of how
to count the number of ways of distributing k balls into n different boxes, under various conditions.
The conditions that are generally imposed are as: (i) The balls can be either different or identical; (ii)
The boxes can be either different or identical; (iii) The distribution can take place either with
exclusion (i.e. no box can contain more than one ball) or without exclusion (i.e. a box may contain
more than one ball). In our discussion the order in which the balls are placed into the boxes is not
important. Now let us discuss the various cases:
Case I: Distributing k different balls into n different boxes, with exclusion, corresponds to
forming a permutation of size k , taken from a set of size n . Therefore, there are n Pk different
ways to distribute k different balls into n different boxes, with exclusion.
Case 2: Distributing k different balls into n different boxes, without exclusion, corresponds to
forming a permutation of size k , with unrestricted repetitions, taken from a set of size n .
Therefore, there are n k different ways to distribute k different balls into n different boxes,
without exclusion.
Case 3: Distributing k identical balls into n different boxes, with exclusion corresponds to
forming a combination of size k , taken from a set of size n . Therefore, there are n C k different
ways to distribute k identical balls into n different boxes, with exclusion.
Case 4: Distributing k identical balls into n different boxes, without exclusion, corresponds to
forming a combination of size k with unrestricted repetitions, taken from a set of size n .
n k 1
Therefore, there are C k different ways to distribute k identical balls into n different boxes,
without exclusion.
Case 5: The number of ways to distribute k different balls into n different boxes, with
exclusion, in such a way that no box is empty, is n! if k n and 0 if k n .
Case 6: The number of ways to distribute k different balls into n different boxes, without
exclusion, in such a way that no box is empty is:
n
C0 ( n 0) k n C1 (n 1) k n C2 (n 2) k ( 1) n 1 ( n Cn )(1) k for k n , and 0 if k n .
Case 7: The number of ways to distribute k identical balls into n different boxes, with
exclusion, in such a way that no box is empty, is 1 if k n and 0 if k n .
Case 8: The number of ways to distribute k identical balls into n different boxes, without
exclusion, in such a way that no box is empty, is k 1 Cn 1 .
Exercise: 6.1
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. How many numbers of five digits can be formed from the numbers 2, 0, 4, 3, 8 when repetition of
digit is not allowed? _____
2. How many numbers can be made with the help of the digits 0, 1, 2, 3, 4, 5 which are greater than
3000 (repetition is not allowed)? _____
3. The number of 5 digit telephone numbers having at least one of their digits repeated is
(a) 90000 (b) 100000 (c) 30240 (d) 69760
4. How many words can be made from the letters of the word ‘COMMITTEE’?
(a) 9! (2!) 2 (b) 9! (2!)3 (c) 9! 2! (d) 9!
5. All the letters of the word ‘EAMCET’ are arranged in all possible ways. The number of such
arrangement in which two vowels are not adjacent to each other is _____.
6. m men and n women are to be seated in a row, so that no two women sit together. If m n , then
the number of ways in which they can be seated is
m !( m 1)! m !( m 1)! ( m 1)!( m 1)!
(a) (b) (c) (d) None of these
( m n 1)! ( m n 1)! ( m n 1)!
7. We are to form different words with the letters of the word ‘INTEGER’. Let m1 be the number of
words in which I and N are never together, and m2 be the number of words which begin with I
and end with R. Then m1 m2 is equal to _____.
8. An n digit number is a positive number with exactly n digits. Nine hundred distinct n digit
numbers are to be formed using only the three digits 2, 5 and 7. The smallest value of n for which
this is possible is _____.
9. The number of numbers that can be formed with the help of the digits 1, 2, 3, 4, 3, 2, 1 so that odd
digits always occupy odd places, is _____.
10. In how many ways can 5 boys and 5 girls sit in a circle so that no boys sit together
(a) (5!) (5!) (b) (4!) (5!) (c) {(5!) (5!)} 2 (d) None of these
11. The number of ways in which 5 beads of different colours form a necklace is _____.
12. The number of ways in which 5 male and 2 female members of a committee can be seated around
a round table so that the two female are not seated together is _____.
13. In a conference of 8 persons, if each person shake hand with the other one only, then the total
number of shake hands shall be _____.
14. To fill 12 vacancies there are 25 candidates of which five are from scheduled caste. If 3 of the
vacancies are reserved for scheduled caste candidates while the rest are open to all, then the
number of ways in which the selection can be made
(a) 5 C3 22 C9 (b) 22 C9 5C3 (c) 22 C3 5C3 (d) None of these
15. There are 10 lamps in a hall. Each one of them can be switched on independently. The number of
ways in which the hall can be illuminated is _____.
16. 10 different letters of English alphabet are given. Out of these letters, words of 5 letters are
formed. How many words are formed when at least one letter is repeated
(a) 99748 (b) 98748 (c) 96747 (d) 97147
17. A man has 10 friends. In how many ways he can invite one or more of them to a party
(a) 10! (b) 2
10 (c) (10!) 1 10
(d) 2 1
18. Numbers greater than 1000 but not greater than 4000 which can be formed with the digits 0, 1, 2,
3, 4 (repetition of digits is allowed), are _____.
19. In the 13 cricket players 4 are bowlers, then how many ways can form a cricket team of 11
players in which at least 2 bowlers included _____.
20. In how many ways can 6 persons to be selected from 4 officers and 8 constables, if at least one
officer is to be included _____.
21. In how many ways can 5 prizes be distributed among four students when every student can take
one or more prizes _____.
22. The number of ways in which 9 persons can be divided into three equal groups is _____.
23. The number of ways dividing 52 cards amongst four players equally, are
(a) (52!) (13!) 4 (b) (52!) {(13!) 2 4!} (c) (52!) {(12!) 4 4!} (d) None of these
24. A question paper is divided into two parts A and B and each part contains 5 questions. The
number of ways in which a candidate can answer 6 questions selecting at least two questions from
each part is _____.
25. There are four balls of different colours and four boxes of colours same as those of the balls. The
number of ways in which the balls, one in each box, could be placed such that a ball doesn't go to
box of its own colour is _____.
26. The number of triangles that can be formed by choosing the vertices from a set of 12 points, seven
of which lie on the same straight line, is _____.
27. A student is allowed to select utmost n books from a collection of (2n 1) books. If the total
number of ways in which he can select one book is 63, then the value of n is _____.
28. The number of divisors of 9600 including 1 and 9600 are _____.
29. A five digit number divisible by 3 has to formed using the numerals 0, 1, 2, 3, 4 and 5 without
repetition. The total number of ways in which this can be done is _____.
30. The number of positive integers which can be formed by using any number of digits from 0, 1, 2,
3, 4, 5 but using each digit not more than once in each number is _____.
31. The number of arrangements of the letters of the word BANANA in which two N's do not appear
adjacently is _____.
32. How many ways are there to arrange the letters in the word GARDEN with the vowels in
alphabetical order? _____
33. 20 persons are invited for a party. In how many different ways can they and the host be seated at a
circular table, if the two particular persons are to be seated on either side of the host?
(a) 20! (b) 2(18!) (c) 18! (d) None of these
34. Five balls of different colours are to be placed in three boxes of different sizes. Each box can hold
all five balls. In how many ways can we place the balls so that no box remains empty? _____
35. A box contains two white balls, three black balls and four red balls. In how many ways can three
balls be drawn from the box if at least one black ball is to be included in the draw? _____
36. A student is to answer 10 out of 13 questions in an examination such that he must choose at least
4 from the first five questions. The number of choices available to him is _____.
37. A man has 7 relatives, 4 women and 3 men. His wife also has 7 relatives, 3 women and 4 men. In
how many ways can they invite 3 women and 3 men so that 3 of them are the man’s relatives and
3 his wife’s _____
38. The sides AB, BC, CA of a triangle ABC have respectively 3, 4 and 5 points lying on them. The
number of triangles that can be constructed using these points as vertices is _____.
39. Six points in a plane be joined in all possible ways by indefinite straight lines, and if no two of
them be coincident or parallel, and no three pass through the same point (with the exception of the
original 6 points). The number of distinct points of intersection is equal to _____.
40. In a plane there are 37 straight lines of which 13 pass through the point A and 11 pass through the
point B. Besides no three lines pass through one point, no line passes through both points A and B
and no two are parallel. Then the number of intersection points the lines have is equal to _____
41. In a polygon no three diagonals are concurrent. If the total number of points of intersection of
diagonals interior to the polygon be 70, then the number of diagonals of the polygon is _____
42. The number of numbers of 4 digits which are not divisible by 5 are
(a) 7200 (b) 3600 (c) 14400 (d) 1800
43. The number of ways in which an examiner can assign 30 marks to 8 questions, awarding not less
than 2 marks to any question is
(a) 21 C16 (b) 30 C16 (c) 21 C7 (d) None of these
44. The number of divisors of the form 4n 2 ( n 0) of the integer 240 is
(a) 4 (b) 8 (c) 10 (d) 3
6.2 Statistics
Statistics may be defined as the science of collection, presentation, analysis and interpretation of
numerical data. It is also defined as the numerical facts or observations collected with definite
purpose. In the language of statistics, one of the most basic concepts is sampling. In most statistical
problems, a specified number of measurements or data or a sample is drawn from a much larger body
of measurements, called the population. The following are the basic terminology used in Statistics:
A variable is a characteristic that changes or varies over time and/or for different individuals or
objects under consideration. There are two types of variables: (i) qualitative (ii) quantitative.
An experimental unit is the individual or object on which a variable is measured. A single
measurement or data value results when a variable is actually measured on an experimental unit.
A population is the set of all measurements of interest to the investigator.
A sample is a subset of measurements selected from the population of interest.
Univariate data result when a single variable is measured on a single experimental unit.
Bivariate data result when two variables are measured on a single experimental unit.
Multivariate data result when more than two variables are measured.
Qualitative variables measure a quality or characteristic on each experimental unit. Qualitative
variables produce data that can be categorized according to similarities or differences in kind; so,
they are often called categorical data. Variables like gender, city, etc. are qualitative variables.
Quantitative variables measure a numerical quantity or amount on each experimental unit. There
are two types of quantitative variable which are given as:
A discrete variable can assume only a finite or countable number of values. Variables such
as number of family members, number of new car sales, and number of defective tires
returned for replacement are all examples of discrete variables.
A continuous variable can assume the infinitely many values corresponding to the points on
a line interval. Variables such as height, weight, time, distance, and volume are continuous
because they can assume values at any point along a line interval. For any two values we pick,
a third value can always be found between them.
measure ‘how often’ in three different ways: (i) The frequency, or number of measurements in each
category; (ii) The relative frequency, or proportion of measurements in each category; (iii) The
percentage of measurements in each category. If n be the total number of measurements in the set,
we can find the relative frequency and percentage using the relationships:
Relative frequency frequency n Percent 100 Relative frequency
The sum of the frequencies is always n ; the sum of the relative frequencies is 1; and the sum of
the percentages is 100%.
Graphs for qualitative data: When the variable is qualitative, the categories should be chosen so
that (i) a measurement will belong to one and only one category (ii) each measurement has a category
to which it can be assigned. For e.g. if we categorize milk according to the type of milk used, we
might use these categories: buffalo, cow, goat, other. Once the measurements have been categorized
and summarized in a statistical table, we can use either a pie chart or a bar chart to display the
distribution of the data.
A pie chart is the familiar circular graph that shows how the measurements are distributed among
the categories. To construct a pie chart, assign one sector of a circle to each category. The angle
of each sector should be proportional to the proportion of measurements (or relative frequency) in
that category. As a circle contains 360°, so Angle Relative frequency 360o to find the angle.
A bar chart shows the same distribution of measurements among the categories, with the height
of the bar measuring how often a particular category was observed.
Graphs for Quantitative Data: Quantitative variables measure an amount or quantity on each
experimental unit. If the variable can take only a finite or countable number of values, it is a discrete
variable. A variable that can assume an infinite number of values corresponding to points on a line
interval is called continuous.
Pie Charts and Bar Charts: Sometimes information is collected for a quantitative variable
measured on different segments of the population, or for different categories of classification. For
example, we might measure the average incomes for people of different age groups, different
genders, or living in different geographic areas of the country. In such cases, we can use pie charts
or bar charts to describe the data, using the amount measured in each category rather than the
frequency of occurrence of each category. The pie chart displays how the total quantity is
distributed among the categories, and the bar chart uses the height of the bar to display the
amount in a particular category.
Line Charts: When a quantitative variable is recorded over time at equally spaced intervals (such
as daily, weekly, monthly, quarterly, or yearly), the data set forms a time series. Time series data
are most effectively presented on a line chart with time as the horizontal axis. The idea is to try to
discern a pattern or trend that will likely continue into the future, and then to use that pattern to
make accurate predictions for the immediate future.
Interpreting Graphs: A distribution is symmetric if the left and right sides of the distribution,
when divided at the middle value, form mirror images.
A distribution is skewed to the right if a greater proportion of the measurements lie to the right
of the peak value. Distributions that are skewed right contain a few unusually large
measurements.
A distribution is skewed to the left if a greater proportion of the measurements lie to the left of
the peak value. Distributions that are skewed left contain a few unusually small measurements.
A distribution is unimodal if it has one peak; a bimodal distribution has two peaks. Bimodal
distributions often represent a mixture of two different populations in the data set.
Relative Frequency Histogram: A relative frequency histogram resembles a bar chart, but it is
used to graph quantitative rather than qualitative data. A relative frequency histogram for a
quantitative data set is a bar graph in which the height of the bar shows “how often” (measured as a
proportion or relative frequency) measurements fall in a particular class or subinterval. The classes or
subintervals are plotted along the horizontal axis. First, divide the interval from the smallest to the
largest measurements into subintervals or classes of equal length. Then find the frequency of each
class and then divide each frequency by total number of classes so we get relative frequency. Based
on relative frequency, draw a bar over each class and we have created a frequency histogram or a
relative frequency histogram, depending on the scale of the vertical axis.
Mean: The arithmetic mean is the central value of the distribution in the sense that positive and
negative deviations from the arithmetic mean balance each other. It is a quantitative average. The
arithmetic mean or average of a set of n measurements is equal to the sum of the measurements
divided by n . To distinguish between the mean for the sample and the mean for the population, we
will use the symbol x for a sample mean and the symbol for the mean of population. The sample
mean x changes from sample to sample but population mean stays the same. Suppose there are n
measurements, x1 , x2 , , xn , on the variable X , then,
n
i 1 xi , n
Mean
n
i1 xi x1 x2 xn (6.3)
Example 6.38 [MN-2009 (2 marks)]: The mean of the cubes of the first n natural numbers is
(a) {n( n 1) 2 } 4 (b) {n ( n 1)(n 2)} 8 (c) ( n4 1) n (d) n3 4
2
13 23 33 n3 1 n( n 1) n( n 1) 2
Solution (a): .
n n 2 4
Median: Median is the central value of the distribution in the sense that the number of values
less than the median is equal to the number of values greater than the median; hence median
is a positional average. The median m of a set of n measurements is the value of x that falls
in the middle position when the measurements are ordered from smallest to largest. If
x1 , x2 , , xn are n values of a variables X , then to find the median, we use the following
algorithm:
Step 1: Arrange the observations x1 , x2 , , xn in ascending or descending order of magnitude
Step 2: Determine the total number of observations, say, n
Step 3: If n is odd, then median is the value of {( n 1) 2}th observation; If n is even, then median is
th
the arithmetic mean of the values of ( n 2)th and ( n 2) 1 observations.
Example 6.40 [MN-2009 (2 marks)]: The following data represent the number of workers suffering
from pneumoconiosis in 10 coal mines.
Mine I II III IV V VI VII VIII IX X
Number 10 16 14 15 14 12 17 13 15 12
The number of mines falling above 50th percentile in terms of workers suffering from pneumoconiosis
is
(a) 2 (b) 3 (c) 4 (d) 5
Solution (c): As the median cut the distribution in half so median is the 50th percentile. Thus in the
given question we have to find how many mines are above the median of suffered workers. As we
have 10 data (which is even), and the data in ascending order as: 10, 12, 12, 13, 14, 14, 15, 15, 16, 17.
th
So the median is the arithmetic mean of the values of (10 2)th and (10 2) 1 data, i.e., arithmetic
mean of 5th and 6th data. So median (14 14) 2 14 . Thus the number of mines having more than
14 workers are 4 (II, IV, VII, IX).
[Similar question was also asked in MT-2010 (1 mark)]
Example 6.41 [CE-2016 (1 mark)]: The spot speeds (expressed in km/hr) observed at a road section
are 66, 62, 45, 79, 32, 51, 56, 60, 53, and 49. The median speed (expressed in km/hr) is _____.
Solution: Arranging the given numbers in an increasing order: 32, 45, 49, 51, 53, 56, 60, 62, 66, 79.
As we have even number of data, i.e. n 10 , so median is arithmetic mean of the values of
( n 2)th (10 2) th 5th observation 53 and {( n 2) 1}th {(10 2) 1}th 6th observation 56 .
Hence median (53 56) 2 54.5 km/hr.
Example 6.42 [MT-2014 (1 mark)]: What is the median value of the following set of numbers? 1, 3,
5, 9, 6, 4, 8 …………….
Solution: As we have 7 data (which is odd) and the data in ascending order as: 1, 3, 4, 5, 6, 8, 9. So
median is {(7 1) 2}th data, i.e., 4th data which is 5.
Mode: The mode or modal value of a distribution is that value of the variable for which the
frequency is maximum. Thus, the mode of a distribution is that value of the variable for which the
frequency variable are clustered densely.
Mode for Discrete frequency distribution: In order to compute the mode of a series of
individual observations, we first convert it into a discrete series frequency distribution by
preparing a frequency table. From the frequency table, we identify the value having maximum
frequency. The value of variable so obtained is the mode or modal value.
Example 6.44: Compute the mode for the following frequency distribution:
Class: 3–6 6–9 9 – 12 12 – 15 15 – 18 18 – 21 21 – 24
Frequency: 2 5 10 23 21 12 3
Solution: We observed that the class 12 – 15 has maximum frequency. Therefore, this is the modal
class. Hence, l 12 , h 3 , f 23 , f1 10 , f 2 21 ; now by using the formula in Step 3, we have
Mode 12 {(23 10) (46 10 21)} 3 14.6 .
In symmetric or normal distribution, as shown in Fig. 6.2(a): Mean Median Mode [This
point was asked in MN-2013 (1 mark)]
In positively skewed distribution, as shown in Fig. 6.2(c): Mode Median Mean [The above
four points were asked in CE-2005 (1 mark)]
Example 6.45 [CH-2013 (1 mark)]: The number of e-mails received on six consecutive days is 11, 9,
18, 18, 4 and 15, respectively. What are the median and the mode for these data?
(a) 18 and 11, respectively (b) 13 and 18, respectively
(c) 13 and 12.5, respectively (d) 12.5 and 18, respectively
Solution (b): As we have total 6 data (which is even) and the data in ascending order is: 4, 9, 11, 15,
th
18, 18. So median is the arithmetic mean of the values of (6 2)th and (6 2) 1 data, i.e., median
is arithmetic mean of 3rd and 4th data, which is (11 15) 2 13 . As the highest frequency in the given
data is for 18 which occurs 2 times; so modal value is 18.
Consider the two distributions shown in Figure 6.3. Both distributions are centred at x 4 , but there
is a big difference in the way the measurements spread out, or vary. The measurements in Figure
6.3(a) vary from 3 to 5; in Figure 6.3(b) the measurements vary from 0 to 8. Measures of variability
can help us to create a mental picture of the spread of the data. We will present some of the more
important ones.
Standard Deviation: One of the most commonly used measures is standard deviation. This value
gives information on how the values of the data set are varying, or deviating, from the mean of the
data set. In other words, the statistical measure of the variability of a distribution around its mean is
referred to as standard deviation [This point was asked in MN-2009, AG-2011 (1 mark)].
Deviations are calculated by subtracting the mean, x , from each of the sample values, x , i.e.
deviation x x . As some values are less than the mean, negative deviations will result, and for
values greater than the mean positive deviations will be obtained. By simply adding the values of the
deviations from the mean, the positive and negative values will cancel to result `in a value of zero. By
squaring each of the deviations, the problem of positive and negative values is avoided. To calculate
the standard deviation, the deviations are squared. These values are summed, divided by the
appropriate number of values and then finally the square root is taken of this result, to counteract the
initial squaring of the deviation. The standard deviation of a population, , of N data items is
N
i 1
( xi )
2
N , where is the population mean (6.5)
N
i 1
( xi ) 2 N N N
x2
i 1 i x
n
i 1 i
2
N2 (6.6)
N N
f x2
i 1 i i n
fx
i 1 i i
2
N 2 , where N i 1 fi
n
(6.7)
Example 6.47 [BT-2016 (1 mark)]: Runs scored by a batsman in five one-day matches are 55, 75,
67, 88 and 15. The standard deviation is _____.
Solution:
xi xi ( xi ) 2 xi 300
55 5 25 xi 5 60 ;
75 15 225
67 7 49 ( xi ) 2 3108
88 28 784
15 45 2025
( xi ) 2 N 3108 5 24.93
Example 6.48 [CS-2011 (2 marks)]: Consider a finite sequence of random values
X x1 , x2 , , xn . Let x be the mean and x be the standard deviation of X . Let another finite
sequence Y of equal length be derived from this as yi a xi b , where a and b are positive
constants. Let y be the mean and y be the selected deviation of this sequence. Which one of the
following sequence is INCORRECT?
(a) Index position of mode of X in X is the same as the index position of mode of Y in Y
(b) Index position of median of X in X is the same as the index position of median of Y in Y
(c) y a x b (d) y a x b
Solution (d): Since mode is that value of the variable for which the frequency is maximum, so if xi at
the k th position is the mode of X then yi at the k th position is the mode Y since yi a xi b , so
option (a) is correct. Also Median is the central value of the distribution in the sense that the number
of values less than the median is equal to the number of values greater than the median so if q th
observation of X is the median then q th observation of Y is also median since yi a xi b , so
n
option (b) is correct. As x i 1 xi n and
n n n
y i 1 yi n i 1 ( axi b) n a i 1 xi n nb n a x b , so option (c) is correct. Since
n n n
x i1 ( xi x )2 n and y i1 ( yi y )2 n i1 (axi b a x b)2 n
n
y a i1 ( xi x )2 n a x ; hence the given option (d) is incorrect.
Variance: The variance of a population of N measurements is the average of the squares of the
deviations of the measurements about their mean . The population variance is denoted by 2 and
is given by the formula 2 N
i 1
( xi )
2
N . Hence the square of standard deviation ( ) is
called as the variance ( 2 ).
Example 6.49 [BT-2012 (2 marks)]: Consider the data set 14, 18, 14, 14, 10, 29, 33, 31, 25. If you
add 20 to each of the values, then
(a) both mean and variance change (b) both mean and variance are unchanged
(c) the mean is unchanged, variance changes (d) the mean changes, the variance is unchanged
Solution (d): As Mean x n
i 1 i
n ; if we add a number ‘ k ’ to each xi , then
n n n
n n
( xi k ) n ( xi ) ( k ) n ( xi ) kn n ( xi ) n k k . Thus
i 1 i 1 i 1 i 1 i 1
adding a constant ‘ k ’ to each data will change the mean by ‘ k ’. Now Variance
2 N
i 1
( xi )
2
2
N ; if we add a number ‘ k ’ to each xi , then N
i 1
( xi k )
2
N;
as ( xi k ) 2 ( xi k k ) 2 ( xi ) 2 , So 2 N
i 1
( xi )
2
2
N , which there is no
change in variance if a constant ‘ k ’ is added to each data.
Statement for Linked Answer Questions 6.50 and 6.51: The abdomen length (in millimetres) was
measured in 15 male fruit flies, and the following data were obtained: 1.9, 2.4, 2.1, 2.0, 2.2, 2.4, 1.7,
1.8, 2.0, 2.0, 2.3, 2.1, 1.6, 2.3, and 2.2.
Example 6.50 [BT-2011 (2 marks)]: Variance (Vx ) for this population of fruit flies as calculated
from the above data shall be
(a) 0.85 (b) 0.25 (c) 0.061 (d) 0.08
Solution (c): Consider the following table with mean (sum of all data) N 31 15 2.067 .
( xi ) Frequency ( fi ) xi2 fi xi f i xi2
1.9 1 3.61 1.9 3.61
2.4 2 5.76 4.8 11.52
2.1 2 4.41 4.2 8.82
2.0 3 4.00 6.0 12
2.2 2 4.84 4.4 9.68
1.7 1 2.89 1.7 2.89
1.8 1 3.24 1.8 3.24
2.3 2 5.29 4.6 10.58
1.6 1 2.56 1.6 2.56
Total N fi 15 xi2 36.6 fi xi 31 fi xi2 64.9
So, Standard deviation, N N
f x2
i 1 i i n
fx
i 1 i i
2
N 2 {15(64.9) 312 } 152 0.25
Example 6.51 [BT-2011 (2 marks)]: The value of standard deviation (SD) will be
(a) 0.061 (b) 0.25 (c) 0.61 (d) 0.85
Solution: From the previous question calculation we have 0.25 .
Coefficient of Variation: Without an understanding of the relative size of the standard deviation
compared to the original data, the standard deviation is somewhat meaningless for use with the
comparison of data sets. To address this problem the coefficient of variation is used. The coefficient
of variation, CV , gives the standard deviation as a percentage of the mean of the data set. The
coefficient of variation for a population is given as,
CV ( ) 100% (6.8)
[Eq. 6.8 was asked in MN-2012 (1 mark)]. The coefficient of variation is often used to compare the
variability of two data sets. It allows comparison regardless of the units of measurement used for each
set of data. The larger the coefficient of variation, the more the data varies.
Example 6.53 [CE-2007 (2 marks)]: If the standard deviation of the spot speed of vehicles in a
highway is 8.8 km/hr and the mean speed of the vehicles is 33 km/hr, the coefficient of variation in
speed is
(a) 0.1517 (b) 0.1867 (c) 0.2666 (d) 0.3646
Solution (c): From the given data we have, standard deviation of the spot speed of vehicles 8.8
km/hr, mean speed of the vehicles 33 km/hr. So from Eq. 6.8, coefficient of variation in speed
CV 8.8 33 0.2666 .
[Similar question was also asked in MN-2010 (1 mark)]
Example 6.54: The results of two tests are shown as: Test 1 (out of 15 marks): 9, 2 ; Test 2
(out of 50 marks): 27, 8 . Compare the variability of these data sets.
2 8
Solution: CVtest 1 100% 100% 22.2% ; CVtest 2 100% 100% 29.6% . Hence,
9 27
the results in the second test show a great variation than those in the first test.
Correlation and Regression: Sometimes, in practical applications, we might come across certain
set of data, where each item of the set may comprise of the values of two or more variables. Suppose
we have a set of 30 students in a class and we want to measure the heights and weights of all the
students. We observe that each individual (unit) of the set assumes two values – one relating to the
height and the other to the weight. Such a distribution in which each individual or unit of the set is
made up of two values is called a bivariate distribution. For e.g. in a class of 60 students the series of
marks obtained in two subjects by all of them. Thus in a bivariate distribution, we are given a set of
pairs of observations, wherein each pair represents the values of two variables.
In a bivariate distribution, we are interested in finding a relationship (if it exists) between the two
variables under study. The concept of ‘correlation’ is a statistical tool which studies the relationship
between two variables and Correlation Analysis involves various methods and techniques used for
studying and measuring the extent of the relationship between the two variables. Two variables are
said to be in correlation if the change in one of the variables results in a change in the other variable.
There are two important types of correlation. They are (1) Positive and Negative correlation and (2)
Linear and Non – Linear correlation.
Positive and Negative Correlation: If the values of the two variables deviate in the same direction
i.e. if an increase (or decrease) in the values of one variable results, on an average, in a corresponding
increase (or decrease) in the values of the other variable the correlation is said to be positive, for e.g.,
household income and expenditure of a family. Correlation between two variables is said to be
negative or inverse if the variables deviate in opposite direction. That is, if the increase in the
variables deviate in opposite direction. That is, if increase (or decrease) in the values of one variable
results on an average, in corresponding decrease (or increase) in the values of other variable, for e.g.,
volume and pressure of perfect gas.
Linear and Non-Linear Correlation: The correlation between two variables is said to be linear if the
change of one unit in one variable result in the corresponding change in the other variable over the
entire range of values. For example consider the following data (2, 7), (4,13), (6,19), (8, 25), (10, 31) ;
thus, for a unit change in the value of x , there is a constant change in the corresponding values of y
and the above data can be expressed by the relation y 3x 1 . In general two variables x and y are
said to be linearly related, if there exists a relationship of the form y a bx , where ‘ a ’ and ‘ b ’ are
real numbers. This is nothing but a straight line when plotted on a graph sheet with different values of
x and y and for constant values of a and b .
The relationship between two variables is said to be non – linear if corresponding to a unit change in
one variable, the other variable does not change at a constant rate. In such cases, if the data is plotted
on a graph sheet we will not get a straight line curve. For example, one may have a relation of the
form y a bx cx 2 or more general polynomial.
Regression Lines (or Regression Models): In case of simple linear regression model (i.e. when there
is only one independent variable and there is linear relationship between the dependent and
independent variable) there are two regression lines as follows:
The product of the two regression coefficients is equal to square of the coefficient of
correlation, i.e. bXY bYX r 2 .
bXY , bYX , r all have the same sign. If r 0 then bXY , bYX are also zero.
The regression lines always intersect at their means.
The angle between the two regression lines depends on the coefficient of correlation, i.e. if
r 0 then the lines are perpendicular to each other; and if r 1 then the regression lines
coincides.
Example 6.55 [XE-2008 (2 marks)]: The two lines of regression of the variables x and y are
4 x 2.4 y 20 and 1.6 x 4 y 12 . The coefficient of correlation between x and y is
(a) 0.49 (b) –0.49 (c) 0.35 (d) –0.35
Solution (b): 4 x 2.4 y 20 ( x 0) 0.6( y 5) which is a regression line x on y , thus
bXY 0.6 and 1.6 x 4 y 12 y 0.4( x 3) which is a regression line y on x , thus bYX 0.4
. So coefficient of correlation, r bXY bYX ( 0.6)( 0.4) 0.489 , we put ve sign because both
bXY and bYX are ve .
Coefficient of Correlation: One of the most widely used statistics is the coefficient of correlation ‘ r
’ which measures the degree of association between the two values of related variables given in the
data set. It takes values from 1 to 1 . If two sets or data have r 1 , they are said to be perfectly
correlated positively if r 1 they are said to be perfectly correlated negatively; and if r 0 they are
uncorrelated. The coefficient of correlation ‘ r ’ is given by the formula
r n ( xy ) x y n x x
2 2
n y y
2 2
.
Regression Analysis: Method of Least Square: If two variables are significantly correlated, and if
there is some theoretical basis for doing so, it is possible to predict values of one variable from the
other. In other words if a regression model is used to express a variable Y as a function of another
variable X then a value of X may be used to estimate a value of Y [This point was asked in ME-
2002 (1 mark)]. This observation leads to a very important concept known as ‘Regression Analysis’.
Regression analysis, in general sense, means the estimation or prediction of the unknown value of one
variable from the known value of the other variable. It is a mathematical measure of the average
relationship between two or more variables in terms of the original units of the data. Suppose we have
a sample of size ‘ n ’ and it has two sets of measures, denoted by x and y . We can predict the values
of ‘ y ’ given the values of ‘ x ’ by using the equation, called the Regression Equation as: y* a bx ,
n ( xy ) x y y b x .
where the coefficients a and b are given by b 2
and a
n x x
2
n
The symbol y * refers to the predicted value of y from a given value of x from the regression
equation. Also the slope of best fit line is ‘ b ’ having intercept ‘ a ’.
Example 6.57 [ME-1998 (2 marks)]: The best fit line using least squares for the data (0, 0), (10, 24),
(20, 36) and (30, 60) is
(a) 2 x y 0 (b) 2 x y 4 0 (c) 2 x y 4 0 (d) None of these
Solution: From the given data we have the following table:
i xi yi xi yi xi2
1 0 0 0 0
2 10 24 240 100
3 20 36 720 400
4 30 60 1800 900
n4 xi 60 yi 120 xi yi 2760 xi2 1400
4(2760) (60)(120) 120 (1.92)(60)
So we have, b 2
1.92 and a 1.2 . Thus the best fit line
4(1400) (60) 4
using the given data is y a bx 1.2 1.92 x . So option (d) is correct.
Example 6.58 [CE-2008 (2 marks)]: Three values of x and y are to be fitted in a straight line in the
form y a bx by the method of least squares. Given x 6 , y 21 , x 2 14 and
2
Coefficient of Determination: The coefficient of determination, r , is useful because it gives the
proportion of the variance of one variable that is predictable from the other variable. It is a measure
that allows us to determine how certain one can be in making predictions from a certain model or
graph. The following points holds for coefficient of determination:
It is the ratio of the explained variation to the total variation
It represents the percent of data that is the closest to the line of best fit.
It is a measures of how well the regression line represent data. If the regression line passes exactly
through every point on the scatter plot, it would be able to explain all of the variation. The further
the line is away from the points, the less it is able to explain.
Exercise: 6.2
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. The mean, median and mode, respectively, for the following ungrouped data: 4, 6, 5, 5, 7, 4, 8, 3,
4, 4, is
(a) 5, 5.5, 4 (b) 5, 4.5, 4 (c) 4.5, 5, 5 (d) 5.5, 5, 5
2. The variance and standard deviation, respectively, of the following ungrouped data: 1, 2, 3, is
(a) 1, 1 (b) 2, 1 (c) 1, 2 (d) 2, 2
3. The following table gives the frequency Number of Order Frequency
distribution of the number of orders 10 – 12 4
received each day during the past 50 days 13 – 15 12
at the office of a mail-order company. 16 – 18 20
Calculate the mean. _____ 19 – 21 14
4. The mode of the frequency distribution table given in Q. 4 is _____.
5. The standard deviation of the frequency distribution table given in Q. 3 is _____.
6. The distribution whose most values are dispersed to the left or right of the mode is classified as
(a) skewed (b) explored (c) bimodal (d) unimodal
Deterministic experiment: Experiments which when repeated under identical conditions produce
the same results or outcomes are known as deterministic experiments. For e.g. when experiments
in science are repeated under identical conditions, we get almost the same result every time.
Random experiment: If an experiment, when repeated under identical conditions, do not produce
the same outcome every time but the outcome in a trial is one of the several possible outcomes
then such an experiment is known as a probabilistic experiment or a random experiment. In a
random experiment, all the outcomes are known in advance but the exact outcome is
unpredictable. For e.g. in tossing of a coin, it is known that either a head or a tail will occur but
one is not sure if a head or a tail will be obtained. So, it is a random experiment.
Event: An event is a subset of a sample space. There are seven types of events:
Simple event: An event containing only a single sample point is called an elementary or
simple event. For e.g. in a single toss of coin, the event of getting a head is a simple event. Here
S {H , T } and E {H } .
Compound events: Events obtained by combining together two or more elementary events
are known as the compound events or decomposable events. For e.g. in a single throw of a
pair of dice the event of getting a doublet, is a compound event because this event occurs if
any one of the elementary events (1,1), (2, 2), (3, 3), (4, 4), (5,5), (6, 6) occurs.
Equally likely events: Events are equally likely if there is no reason for an event to occur in
preference to any other event. For e.g., if an unbiased die is rolled, then each outcome is
equally likely to happen, i.e. all elementary events are equally likely.
Mutually exclusive or disjoint events: Events are said to be mutually exclusive or disjoint if
the occurrence of any one of them prevents the occurrence of all the others. For e.g., E
getting an even number, F getting an odd number, these two events are mutually
exclusive, because, if E occurs we say that the number obtained is even and so it cannot be
odd i.e., F does not occur. A1 and A2 are mutually exclusive events if A1 A2 .
Mutually non-exclusive events: The events which are not mutually exclusive are known as
compatible events or mutually non – exclusive events.
Independent events: Events are said to be independent if the happening (or non-happening)
of one event is not affected by the happening (or non-happening) of others. For e.g. if two
dice are thrown together, then getting an even number on first is independent to getting an
odd number on the second.
Dependent events: Two or more events are said to be dependent if the happening of one
event affects (partially or totally) other event. For e.g. suppose a bag contains 5 white and 4
black balls. Two balls are drawn one by one. Then two events that the first ball is white and
second ball is black are independent if the first ball is replaced before drawing the second
ball. If the first ball is not replaced then these two events will be dependent because second
draw will have only 8 exhaustive cases.
Exhaustive number of cases: The total number of possible outcomes of a random experiment in
a trial is known as the exhaustive number of cases. For e.g., in throwing a die the exhaustive
number of cases is 6, as any one of the six faces marked with 1, 2, 3, 4, 5, 6 may come uppermost.
Favourable number of cases: The number of cases favourable to an event in a trial is the total
number of elementary events such that the occurrence of any one of them ensures the happening
of the event. For e.g., in drawing two cards from a pack of 52 cards, the number of cases favourable to
drawing 2 queens is 4 C2 .
Mutually exclusive and exhaustive system of events: Let S be the sample space associated
with a random experiment. Let E1 , E2 , , En are elementary events associated with a random
experiment such that (i) Ei E j for i j ; (ii) E1 E2 En S . Then the collection of
events E1 , E2 , , En is said to form a mutually exclusive and exhaustive system of events. Hence, the
collection of elementary events associated with a random experiment always form a system of
mutually exclusive and exhaustive system of events. In this system,
P ( E1 E2 En ) P ( E1 ) P ( E2 ) P ( En ) 1 .
Independent events are always taken from different experiments, while mutually exclusive events
are taken from a single experiment.
Independent events can happen together while mutually exclusive events cannot happen together.
Independent events are connected by the word “and” but mutually exclusive events are connected
by the word “or”.
In other words, if a random experiment results in n mutually exclusive, equally likely and exhaustive
outcomes, out of which m are favourable to the occurrence of an event A , then the probability of
m Number of outcomes favourable to A
occurrence of A is given by P ( A) .
n Number of total outcomes
Further, if A denotes negative of A i.e., event that A does not happen, then for above cases m , n ;
we have,
P ( A ) ( n m ) n 1 ( m n ) 1 P ( A) P ( A) P ( A ) 1 (6.10)
Some important remarks about Coins, Dice, Playing cards and Envelopes
Coins: A coin has a head side and a tail side. If an experiment consists of more than a coin, then
coins are considered to be distinct if not otherwise stated. Number of exhaustive cases of tossing
n coins simultaneously (or of tossing a coin n times) 2 n .
Dice: A die (cubical) has six faces marked 1, 2, 3, 4, 5, 6 . We may have tetrahedral (having four
faces 1, 2, 3, 4 ) or pentagonal (having five faces 1, 2, 3, 4, 5 ) die. As in the case of coins, if we have
more than one die, then all dice are considered to be distinct if not otherwise stated. Number of
exhaustive cases of throwing n dice simultaneously (or throwing one dice n times) 6 n .
Playing cards: A pack of playing cards usually has 52 cards. There are 4 suits (Spade, Heart,
Diamond and Club) each having 13 cards. There are two colours red (Heart and Diamond) and
black (Spade and Club) each having 26 cards. In thirteen cards of each suit, there are 3 face cards
or coart cards namely king, queen and jack. So there are in all 12 face cards (4 kings, 4 queens
and 4 jacks). Also there are 16 honour cards, 4 of each suit namely ace, king, queen and jack.
Theorems on Probability:
When events are not mutually exclusive: If A and B are two events which are not mutually
exclusive, then P ( A B ) P ( A) P ( B ) P ( A B ) .
P ( A B ) P ( A) P ( B ) [This point was asked in CS-1994 (1 mark)]
For any three events A, B, C
P ( A B C ) P( A) P ( B ) P (C ) P ( A B ) P ( B C ) P (C A) P ( A B C )
When events are mutually exclusive: If A and B are mutually exclusive events, then
n( A B ) 0 P( A B ) 0 . Hence P ( A B) P ( A) P( B ) .
For any three events A, B, C which are mutually exclusive,
P ( A B ) P ( B C ) P (C A) P( A B C ) 0
P( A B C ) P ( A) P( B ) P (C ) .
The probability of happening of any one of several mutually exclusive events is equal to the
sum of their probabilities, i.e., if A1 , A2 An are mutually exclusive events, then
P ( A1 A2 An ) P ( A1 ) P( A2 ) P ( An ) i.e., P Ai P ( Ai ) .
When events occur with random experiment: Let A and B be two events associated with a
random experiment, then
P ( A B ) P ( B) P ( A B) ; and P ( A B ) P ( A) P ( A B)
If B A , then P ( A B ) P ( A) P ( B) ; and P ( B ) P ( A)
Similarly if A B , then ( A B ) P ( B ) P ( A) ; and P ( A) P ( B )
Probability of occurrence of neither A nor B is P ( A B ) P ( A B ) 1 P ( A B ) .
( 1) n 1 P ( A1 A2 An ) .
P A
n
i 1 i
n
i 1
P ( Ai ) ( n 1) P A
n
i 1 i
n
i 1
P( Ai )
Example 6.59 [CS-1994 (1 mark)]: Let A , B and C be independent events which occur with
probabilities 0.8, 0.5 and 0.3 respectively. The probability of occurrence of at least one of the event is
………..
Solution: The probability of occurrence of at least one of the event is
P ( A B C ) P( A) P ( B ) P (C ) P ( A B ) P ( B C ) P ( A C ) P ( A B C ) . As all the
event are independent so P ( A B C ) P( A) P ( B ) P (C ) (0.8)(0.5)(0.3) 0.12 and
P ( A B ) P ( A) P ( B) (0.8)(0.5) 0.4 , similarly P( B C ) 0.15 and P ( A C ) 0.21 . Thus
P ( A B C ) 0.8 0.5 0.3 0.4 0.15 0.21 0.12 0.96 .
Example 6.60 [CS-1994 (1 mark)]: The probability of an event B is P1 . The probability that events
A and B occur together is P2 while the probability that A and B occur together is P3 . The
probability of the event A in terms of P1 , P2 and P3 is ……….
Solution: From the given data we have, P ( B) P1 , P ( A B ) P2 , P ( A B ) P3 . As for a random
experiment, P ( A B ) P ( A) P ( A B ) P3 P( A) P2 P ( A) P2 P3 .
Example 6.61 [CS-1995 (1 mark)]: The probability that a number selected at random between 100
and 999 (both inclusive) will not contain the digit 7 is
(a) 16 25 (b) (9 10)3 (c) 27 75 (d) 18 25
Solution (d): As there are 900 numbers between 100 and 999. Now there are three digits, lets us find
the number of digits with no 7; the first digit can be chosen from 8 digits, either 1 2 3 4 5 6 8 or 9 (we
cannot choose 0), the second digit can be chosen from 9 digits (0 to 9 excluding 7), the third digit can
be chosen from any of the 9 digits (excluding 7) so the total number of possibilities or numbers is
8 9 9 648 . So probability of selecting at random a number with no digit 7 is 648 900 18 25 .
Example 6.62 [CS-1996 (2 marks)]: The probability that top and bottom cards of a randomly
shuffled deck are both aces is
(a) (4 52) (4 52) (b) (4 52) (3 52) (c) (4 52) (3 51) (d) (4 52) (4 51)
Solution: Probability that top card of a randomly shuffled deck is an ace is 4 52 . Now we have 51
cards and 3 aces so the probability that bottom card is an ace is 3 51 . As both events are independent
so probability that top and bottom cards of a randomly shuffled deck are both aces is (4 52) (3 51) .
Example 6.63 [CS-1996 (1 mark)]: Two dice are thrown simultaneously. The probability that at least
one of them will have 6 facing up is
(a) 1 36 (b) 1 3 (c) 25 36 (d) 11 36
Solution (d): Two dice are thrown simultaneously so total number of sample space is n( S ) 62 36 .
Now let E be the event of getting at least one of them will have 6 facing up so the elements of
E :{(1, 6), (2, 6), , (5, 6), (6,1), (6, 2), (6, 6)} so n( E ) 11 and thus P ( E ) n( E ) n( S ) 11 36 .
Example 6.64 [CS-1997 (1 mark)]: The probability that it will rain today is 0.5. The probability that
it will rain tomorrow is 0.6. The probability that it will rain either today or tomorrow is 0.7. What is
the probability that it will rain today and tomorrow?
(a) 0.3 (b) 0.25 (c) 0.35 (d) 0.40
Solution (d): Let A be the event that it will rain today, and B be the event that it will rain tomorrow.
So from the given data we have P ( A) 0.5 , P ( B) 0.6 and P( A B ) 0.7 . We have to find
P ( A B ) P ( A) P( B ) P ( A B ) 0.5 0.6 0.7 0.4 .
Example 6.65 [CS-1998 (1 mark)]: A die is rolled three times. The probability that exactly one odd
number turns up among the three outcomes is
(a) 1/6 (b) 3/8 (c) 1/8 (d) ½
Solution (b): Let O be the event of outcome of an odd number; and E be the event of outcome of an
even number. When a die is rolled the outcome is either 3 odd numbers or 3 even numbers; so
P (O) P ( E ) 3 6 1 2 . Now in three tries exactly one odd number turns up among the three
outcomes as: {OEE , EOE , EEO} so probability that exactly one odd number turns up among the three
111 111 111 3
outcomes is P (O ) P ( E ) P ( E ) P ( E ) P (O) P ( E ) P ( E ) P ( E ) P (O ) .
222 222 222 8
Example 6.66 [ME-1998 (1 mark)]: The probability that two friends share the same birth-month is
(a) 1 6 (b) 1 12 (c) 1 144 (d) 1 24
Solution (b): Any person can born in any of the 12 months of the year, so total number of possible
combinations is 12 12 . This set contains 12 pairs of identical months, for e.g., (January, January),
which are our favourable cases. Thus the required probability is 12 (12 12) 1 12 .
Example 6.67 [CS-2000 (1 mark)]: The minimum number of cards to be dealt from an arbitrarily
shuffled deck of 52 cards to guarantee that three cards are from same suit is
Example 6.68 [CS-2001 (2 marks)]: Seven (distinct) car accidents occurred in a week. What is the
probability that they all occurred on the same day?
1 1 1 7
(a) 7 (b) 6 (c) 7 (d) 7
7 7 2 2
Solution (b): Probability of 1 accident in a 7 days is 1 7 . So probability of 7 accidents in 7 days
considering 1 accident in a day is (1 7)(1 7) (7 times) 1 77 . Now consider all 7 accidents in a day
so its probability is 7(1 77 ) 1 76 .
Example 6.69 [ME-2001 (2 marks)]: An unbiased coin is tossed three times. The probability that the
head turns up in exactly two cases is
(a) 1 9 (b) 1 8 (c) 2 3 (d) 3 8
Solution (d): As unbiased coin is tossed three times, so we have 23 8 number of sample space. Out
of 8 events, we have exactly two heads as {HHT , HTH , THH } . So the required probability is 3 8 .
Example 6.70 [CS-2002 (2 marks)]: Four fair coins are tossed simultaneously. The probability that
at least one head and one tail turn up is
1 1 7 15
(a) (b) (c) (d)
16 8 8 16
4
Solution (c): As when four coins are tossed the total number of cases will be 2 16 . Now among 16
possibilities the condition at least one head and one tail is not met by two possibilities, i.e.,
{HHHH , TTTT } ; so total number of favourable cases for at least one head and one tail will be
16 2 14 . Thus the required probability will be 14 16 7 8 .
Example 6.71 [CE-2003 (1 mark)]: A box contains 10 screws, 3 of which are defective. Two screws
are drawn at random with replacement. Probability that none of the two screws is defective will be
(a) 100% (b) 50% (c) 49% (d) None of these
Solution (c): Let A and B be the events of drawing first and second, respectively, non-defective
screws. So for event A , P ( A) 7 10 ; also we are sampling with replacement so the situation before
the second drawing is same as at the beginning and thus P ( B ) 7 10 . As both the events are
independent so P ( A B ) P( A) P ( B ) 49 100 , which is the required probability.
Example 6.72 [ME-2003 (2 marks)]: A box contains 5 black and 5 red balls. Two balls are randomly
picked one after another from the box, without replacement. The probability for both balls being red is
(a) 1 90 (b) 1 5 (c) 19 90 (d) 2 9
Solution (d): The probability of drawing the first red ball is 5 10 . As our case is without replacement,
so we now have 5 black and 4 red ball; thus the probability of drawing the 2nd red ball is 4 9 . So the
required probability is (5 10)(4 9) 2 9 .
[Similar question was also asked in ME-2006 (1 mark)]
Example 6.73 [ME-2004 (2 marks)]: From a pack of regular playing cards, two cards are drawn at
random. What is the probability that both cards will be Kings, if the first card is NOT replaced?
(a) 1 26 (b) 1 52 (c) 1 169 (d) 1 221
Solution (d): As in pack of 52 cars we have 4 kings. So the probability of drawing first king is 4 52 .
As our case is without replacement, so we now have 3 kings and total 51 cards; so the probability of
drawing the second king is 3 51 . Thus the required probability is (4 52)(3 51) 1 221 .
Example 6.74 [EC-2005 (1 mark)]: A fair dice is rolled twice. The probability that an odd number
will follow an even number is
(a) 1 2 (b) 1 6 (c) 1 3 (d) 1 4
Solution (d): A dice is rolled two times so our sample space is 62 36 . Let E be the event for which
an odd number will follow an even number; so the elements of E are {(2,1), (2, 3), (2, 5), (4,1),
(4, 3), (4, 5), (6,1), (6,3), (6,5)} n( E ) 9 ; thus the required probability is 9 36 1 4 .
Example 6.75 [EE-2005 (2 marks)]: A fair coin is tossed three times in succession. If the first toss
produces a head, then the probability of getting exactly two heads in three tosses is:
(a) 1 8 (b) 1 2 (c) 3 8 (d) 3 4
Solution (b): It is given that on tossing a coin for the first time the outcome is head. Now the
question: getting exactly two head in three tosses, is same as getting exactly one head in two tosses.
So our sample space is S :{HH , HT , TH , TT } n( S ) 4 and if E is the event of getting exactly one
head in two tosses then E :{HT , TH } n( E ) 2 and thus P ( E ) 2 4 1 2 which is the probability
of getting exactly two heads in three tosses given that first toss is head.
Example 6.76 [ME-2005 (2 marks)]: A single die is thrown twice. What is the probability that the
sum is neither 8 nor 9?
(a) 1 9 (b) 5 36 (c) 1 4 (d) 3 4
Solution: The total number of sample space when a die is thrown twice is 62 36 . Let E be the
event of getting the sum of numbers on the dice is 8 or 9, then the elements of E are
{(2, 6), (3,5), (4, 4), (5,3), (6, 2), (3, 6), (4,5), (5, 4), (6,3)} n( E ) 9 . So the probability of getting 8
and 9 is 9 36 1 4 . Thus the probability that the sum is neither 8 nor 9 is 1 1 4 3 4 .
Example 6.77 [CE-2006 (2 marks)]: There are 25 calculators in a box. Two of them are defective.
Suppose 5 calculators are randomly picked for inspection (i.e., each has the same chance of being
selected), what is the probability that only one of the defective calculators will be included in the
inspection?
(a) 1 2 (b) 1 3 (c) 1 4 (d) 1 5
25
Solution: As total number of ways for picking 5 out of 25 calculators is C 2 . Now we have to pick 4
23
out of 23 non-defective calculator, which can be done in C4 ; also we have to pick 1 out of 2
defective calculators which can be done in 2 C1 ; so total number of ways in which 5 calculators (4
23
non-defective and 1 defective) can be picked is C4 2 C1 ways. Thus required probability is
( 23 C4 2 C1 ) ( 25 C5 ) 1 3 .
[Similar question was also asked in CS-1995 (2 marks)]
Example 6.78 [CS-2007 (2 marks)]: Suppose we uniformly and randomly select a permutation from
the 20! permutations of 1, 2, 3, …, 20. What is the probability that 2 appears at an earlier position
than any other even number in the selected permutation?
(a) 1 2 (b) 1 10 (c) (9!) (20!) (d) None of these
Solution (b): Number of permutation with ‘2’ in the 1st position is 19! 10 P0 (19!) (fill the first place
with ‘2’ and 19 places after ‘2’ with 19 remaining numbers). Number of permutation with ‘2’ in the
2nd position is 10 18! 10 P1 (18!) (fill the first place with any of the 10 odd numbers and the 18 places
after the 2 with 18 of the remaining numbers in 18! ways). Similarly, number of permutation with ‘2’
in 3rd position is 10 9 17! 10 P2 (17!) (fill the first 2 places with 2 of the 10 odd numbers and then
the remaining 17 places with remaining 17 numbers in 17!). We repeat this process until ‘2’ is in the
11th place. After that it is not possible to satisfy the given condition as there are only 10 odd numbers
available to fill before the ‘2’. So the desired number of permutation which satisfy the given condition
is 10 P0 (19!) 10 P1 (18!) 10 P9 (10!) 10 P10 (9!) . As total number of ways in which these 20 numbers
20
are arrange is P20 20! . So the probability of happening of the given condition is
10
P0 (19!) 10 P1 (18!) 10 P9 (10!) 10 P10 (9!) 1
20
. Alternate way: The odd numbers do not matter
P20 10
here. The probability that 2 comes before the other 9 evens is
(number of ways to pick 2)(number of ways to pick remaining evens) (1)(9!) 1
.
(number of ways to order 10 evens) 10! 10
Example 6.79 [PI-2007 (1 mark)]: Two cards are drawn at random in succession, with replacement,
from a deck of 52 well shuffled cards. Probability of getting both ‘Aces’ is
(a) 1/169 (b) 2/169 (c) 1/13 (d) 2/13
Solution (a): As there are four aces in a deck of 52 cards and we have the case of replacement. So the
probability of getting both aces is (4 52) (4 52) 1 169 .
Example 6.80 [CE-2008 (2 marks)]: A person on a trip has a choice between private car and public
transport. The probability of using a private car is 0.45. While using the public transport, further
choices available are bus and metro, out of which the probability of commuting by a bus is 0.55. In
such a situation, the probability (rounded up to two decimals) of using a car, bus and metro,
respectively would be
(a) 0.45, 0.30 and 0.25 (b) 0.45, 0.25 and 0.30
(c) 0.45, 0.55 and 0.00 (d) 0.45, 0.35 and 0.20
Solution (a): Let C be car; B be bus; and M be metro. It is given that P (C ) {P( B ) P ( M )} 1 ,
P (C ) 0.45 P ( B ) P ( M ) 1 P (C ) 0.55 . Now while using the public transport the probability
of commuting by a bus is 0.55, which means that P ( B ) 0.55 {P ( B ) P ( M )} 0.55 0.55 0.30 .
So P ( M ) 0.55 P( B) 0.25 .
Example 6.81 [CS-2008 (2 marks)]: Aishwarya studies either Computer Science or Mathematics
everyday. If she studies Computer Science on a day, then the probability that the studies Mathematics
the next day is 0.6. If she studies Mathematics on a day, then the probability that the studies Computer
Science the next day is 0.4. Given that Aishwarya studies Computer Science on Monday, what is the
probability that she studies Computer Science on Wednesday?
(a) 0.24 (b) 0.36 (c) 0.40 (d) 0.60
Solution (c): Aishwarya studies computer science on Monday, so the probability that she studies
mathematics on Tuesday is 0.6 and the probability that she studies computer science on Tuesday is
1 0.6 0.4 . The probability that she studies mathematics on Tuesday and computer science on
Wednesday is 0.6 0.4 0.24 . Now the probability that she studies computer science on Tuesday and
computer science on Wednesday is 0.4 0.4 0.16 . So the required probability that she studies
computer science on Wednesday is 0.24 0.16 0.40 .
Example 6.82 [ME-2008, TF-2009 (1 mark)]: A coin is tossed 4 times. What is the probability of
getting heads exactly 3 times?
Example 6.83 [PI-2008 (2 marks)]: In a game, two players X and Y toss a coin alternatively.
Whosoever gets a ‘head’ first, wins the game and the game is terminated. Assuming that player X
starts the game, the probability of player X winning the game is
(a) 1/3 (b) 1/2 (c) 2/3 (d) ¾
Solution (c): Case 1: X starts the game and win, so P1 1 2 . Case 2: X starts the game and loose,
then Y tosses and coin shows tail, then X tosses and coin shows head thus wins, so
P2 (1 2)(1 2)(1 2) 1 23 . Case 3: X starts the game and loose, then Y tosses and coin shows tail,
then X tosses and coin shows tail, then Y tosses and coin shows tail, then X tosses and coin shows
head thus wins, so P3 (1 2)(1 2)(1 2)(1 2)(1 2) 1 25 . And so one …. So the required probability is
P P1 P2 P3 1 2 1 2 3 1 25 (1 2) {1 1 2 2 } 2 3 .
Example 6.84 [TF-2008 (2 marks)]: Two dices are thrown simultaneously. The probability that the
total number of dots is equal to 4 is
(a) 1 6 (b) 1 12 (c) 1 18 (d) 1 36
Solution (b): When two dice are thrown then we have 62 36 number of sample space. Let E be
the event of getting sum of the numbers on two dice as 4. So the elements of E are
{(1, 3), (2, 2), (3,1)} n( E ) 3 . Thus P ( E ) 3 36 1 12 .
[Similar questions were also asked in ME-2002, TF-2014 (1 mark)]
Example 6.85 [CS-2009 (2 marks)]: An unbalanced dice (with 6 faces, numbered from 1 to 6) is
thrown. The probability that the face value is odd is 90% of the probability that the face value is even.
The probability of getting any even numbered face is the same. If the probability that the face is even
given that it is greater than 3 is 0.75, which one of the following options is closest to the probability
that the face value exceeds 3?
(a) 0.453 (b) 0.468 (c) 0.485 (d) 0.492
Solution (b): Let E and O are the events of getting an even and odd numbers, respectively. So we
have P ( E ) x , P (O ) 90% of x 0.9 x ; but P ( E ) P (O ) 1 x 0.9 x 1 x 10 19 . As we
have three even numbers and probability of getting an even number is same. So
P (2) P (4) P(6) (10 19)(1 3) 10 57 . Now
P ( E 3 and number 3) P ( E 3) P(number 3) 0.75 {P(4) P (6)} P (number 3)
P(number 3) 0.75 (10 57 10 57) 0.75 (20 57) 0.468 .
Example 6.86 [EC-2009 (1 mark)]: A fair coin is tossed 10 times. What is the probability that
ONLY the first two tosses will yield heads?
(a) (1 2) 2 (b) 10 C2 (1 2) 2 (c) (1 2)10 (d) 10 C2 (1 2)10
Solution (c): As a fair coin is tossed 10 times so total number of sample space is n( S ) 210 . Let E
be the event of getting first two tosses as ‘Heads’ and other eight as ‘Tail’; this will yields as
{H , H , T , T , T , T , T , T , T , T } so n( E ) 1 . Thus P ( E ) 1 210 .
Example 6.87 [EE-2009 (2 marks)]: Assume for simplicity that N people, all born in April (a
month of 30 days), are collected in a room. Consider the event of at least two people in the room
being on the same date of the month, even if in different years, e.g. 1980 and 1985. What is the
smallest N so that the probability of this event exceeds 0.5?
Example 6.88 [TF-2009 (1 mark)]: Probability of getting 16 in one throw with 3 dice is
(a) 1 8 (b) 3 16 (c) 1 36 (d) 1 108
Solution: When three dice are thrown then we have 63 216 number of sample space. Let E be the
event of getting sum of the numbers on three dice as 16. So the elements of E are
{(4, 6, 6), (6, 4, 6), (6, 6, 4), (5,5, 6), (5, 6,5), (6, 5,5)} n( E ) 6 . Thus P ( E ) 6 216 1 36 .
Example 6.89 [CE-2010 (1 mark)]: Two coins are simultaneously tossed. The probability of two
heads simultaneously appearing is
(a) 1 8 (b) 1 6 (c) 1 4 (d) 1 2
Solution: Our sample space is S :{HH , HT , TH , TT } , so n( S ) 4 ; if E be the even of occurring
two heads then n( E ) 1 . Thus P ( E ) n( E ) n( S ) 1 4 .
Example 6.90 [CS-2010 (2 marks)]: What is the probability that divisor of 1099 is a multiple of 1096
?
(a) 1 625 (b) 4 625 (c) 12 625 (d) 16 625
Solution (a): Divisors of 1099 are of the form 2a 5b , where a and b vary from 0 to 99 each, so
there are 10000 divisors of 1099 . Now any of those divisors will be a multiple of 1096 if both a and
b are at least 96, i.e., 96 or 97 or 98 or 99. So both a and b have 4 choices each, and so there are
4 4 16 divisors which are multiple of 1096 . So the required probability is 16 10000 1 625 .
Example 6.91 [CH-2010 (2 marks)]: A box contains three red and two black balls. Four balls are
removed from the box one by one without replacement. The probability of the ball remaining in the
box being red, is
(a) 609 625 (b) 3 5 (c) 2 5 (d) 81 625
Solution (b): Let E be the event for red ball remaining; and S be the event for only 1 ball
remaining. So number of ways in which event E occurs is when 2 red balls and 2 black balls
selected, i.e., n ( E ) 3C2 2 C2 ; also number of ways in which event S occurs is when 4 out of 5
balls selected, i.e., n ( S ) 5C4 , which is our sample space. So the required probability is
P ( E ) {n( E )} {n( S )} 3
C2 2 C2 C 3 5.
5
4
Example 6.92 [ME-2010 (2 marks)]: A box contains 2 washers, 3 nuts and 4 bolts. Items are drawn
from the box at random one at a time without replacement. The probability of drawing 2 washers first
followed by 3 nuts and subsequently the 4 bolts is
Example 6.93 [MN-2010 (2 marks)]: The probabilities of hitting a target by A and B are 1/3 and 2/5,
respectively. A shoots at the target once, followed by B shooting at the target once. The probability of
hitting the target is
(a) 2 15 (b) 5 15 (c) 8 15 (d) 9 15
Solution (a): Let E , F are the events of hitting a target by A and B , respectively. As both the
events are independent so P( E and F ) P ( E F ) (1 3)(2 5) 2 15 .
Example 6.94 [CE-2011 (1 mark)]: There are two containers, with one containing 4 Red and 3
Green balls and the other containing 3 Blue and 4 Green balls. One ball is drawn at random from each
container. The probability that one of the balls is Red and the other is Blue will be
(a) 1 7 (b) 9 49 (c) 12 49 (d) 3 7
Solution (c): Let R and B be the events of drawing red and blue balls, respectively. So we have to
pick red ball from first container so P ( R ) 4 7 ; and blue ball from second container so P ( B ) 3 7 .
As both the events are independent so P ( R and B ) P ( R B ) P ( R ) P( B ) (3 7)(4 7) 21 49 .
Example 6.95 [CS-2011 (1 mark)]: If two fair coins are flipped and at least one of the outcomes is
known to be a head, what is the probability that both outcomes are head?
(a) 1 3 (b) 1 4 (c) 1 2 (d) 2 3
Solution (a): As our sample space is {HH , HT , TH } ; so the probability that both outcomes are head
is 1 3 .
Example 6.96 [CS-2011 (2 marks)]: A deck of 5 cards (each carrying a distinct number from 1 to 5)
is shuffled thoroughly. Two cards are then removed one at a time from the deck. What is the
probability that the two cards are selected with the number on the first card being one higher than the
number on the second card?
(a) 1 5 (b) 4 25 (c) 1 4 (d) 2 5
Solution (a): Total number of ways in which 2 from 5 cards can be selected is 5 P2 20 ways. Now
let E be the event of selecting two cards s.t. number on the 1st card being one higher than number on
the 2nd card, then E :{(2,1), (3, 2), (4,3), (5, 4)} ; so P ( E ) 4 20 1 5 be the required probability.
Example 6.97 [EC-2011 (1 mark)]: A fair dice is tossed two times. The probability that the second
toss results in a value that is higher than the first toss is
(a) 2 36 (b) 2 6 (c) 5 12 (d) 1 2
Solution (c): As a dice is rolled two times so our sample space is 62 36 . Let E be the event for
which the second toss results in a value that is higher than the first toss; so the elements of E are
{(2,1), (3,1), (3, 2), (4,1), (4, 2), (4,3), (5,1), , (5, 4), (6,1), , (6,5)} , where the first number
corresponds to the second toss and second number corresponds to the first toss. So
n( E ) 1 2 4 4 5 15 . Thus the required probability is P ( E ) n( E ) n( S ) 15 36 5 12 .
Example 6.98 [IN-2011 (2 marks)]: The box 1 contains chips numbered 3, 6, 9, 12, and 15. The box
2 contains chips numbered 6, 11, 16, 21 and 26. Two chips, one from each box are drawn at random.
The numbers written on these chips are multiplied. The probability for the product to be an even
number is
(a) 6 25 (b) 2 5 (c) 3 5 (d) 19 25
Solution (d): Let x be the product of the numbers from box 1 and 2. For x to be even, the numbers
from both the boxes should not turn out to be odd simultaneously. So the required probability is
P ( x is even) 1 P ( x is odd) 1 (3 5) (2 5) 19 25 .
Example 6.99 [ME-2011 (2 marks)]: An unbiased coin is tossed five times. The outcome of each
toss is either a head or a tail. The probability of getting at least one head is
(a) 1 32 (b) 13 32 (c) 16 32 (d) 31 32
Solution (d): Total number of sample space when a coins is tossed 5 times is 25 32 . Let E be the
event of getting at least 1 head; and E1 be the event of getting no head. So the elements of E1 is
{TTTTT } n( E1 ) 1 . So the required probability is P ( E ) 1 P( E1 ) 1 1 32 31 32 .
[Similar question was also asked in ME-2009 (1 mark)]
Example 6.100 [MT-2011 (2 marks)]: A box contains 5 white balls and 3 red balls. Two balls are
withdrawn from the box randomly, one after another (without replacement). The probability that the
two balls withdrawn are of different colour is
(a) 15/64 (b) 25/64 (c) 25/56 (d) 30/56
8
Solution (d): As 2 out of 8 balls are selected in C 2 ways. Let E be the event of selecting 1 out of 5
white balls and 1 out of 3 red balls. So number of ways in which the event E occurs is
n ( E ) 5C1 3C1 . Thus the required probability is P ( E ) 5
C1 3C1 C 15 28 30 56 .
8
2
Example 6.101 [TF-2011 (1 mark)]: Probability of occurrence of two events E1 and E2 is 0.25 and
0.5, respectively. The probability of their simultaneous occurrence is 0.14. The probability that neither
E1 nor E2 occurs is
(a) 0.11 (b) 0.25 (c) 0.39 (d) 0.86
Solution (c): P ( E1 ) 0.25 and P ( E2 ) 0.5 , P ( E1 E2 ) 0.14 ; we have to find
P ( E1c E2c ) 1 P ( A B ) 1 P ( A) P ( B ) P ( A B ) 1 0.25 0.5 0.14 0.39 .
Example 6.102 [TF-2011 (2 marks)]: A garment factory manufactures shirts. From the past history,
it is known that 8 out of 100 collars and 5 out of 100 sleeves are defective. The probability that the
assembled shirt will NOT have either of these defects is ……
Solution: Let E and F are the events of defective collars and sleeves, respectively. So
P ( E ) 8 100 0.08 and P ( F ) 5 100 0.05 . As both the events are independent so
P ( E F ) P ( E ) P ( F ) (0.08)(0.05) 0.004 . Thus a shirt having collar or sleeve to be defective is
P ( E F ) P( E ) P ( F ) P ( E F ) 0.08 0.05 (0.08)(0.05) 0.126 . Hence the probability that
no shirt is defective is 1 P ( E F ) 1 0.126 0.874 .
Example 6.103 [ME-2012, PI-2012 (2 marks)]: A box contains 4 red balls and 6 black balls. Three
balls are selected randomly from the box one after another without replacement. The probability that
the selected set contains one red ball and two black balls is
(a) 1 20 (b) 1 12 (c) 3 10 (d) 1 2
Solution (b): The selection of one red and two black balls balls is done as: BBR , BRB , RBB . So
probability of the event BBR is (6 10) (5 9) (4 8) 120 720 1 6 . Similarly probability of BRB
and RBB is 1 6 and 1 6 , respectively. So total probability of getting two black balls and one red ball
is (1 6) (1 6) (1 6) (3 6) 1 2 .
[Similar question was also asked in ME-1997 (2marks)]
Example 6.104 [CS-2012 (2 marks)]: Suppose a fair six-sided die is rolled once. If the value on the
die is 1, 2, or 3, the die is rolled a second time. What is the probability that the sum total of values that
turn up is at least 6?
(a) 10 21 (b) 5 12 (c) 2 3 (d) 1 6
Solution (b): The situation given in the problem can be drawn as in the
following figure. When die is rolled then we have 6 outcomes. If the outcome is
either 1 or 2 or 3 then the die will again rolled. So if in the first roll the outcome
is 1 then in the second roll the outcome can be any one from 1 to 6; but the sum
of the outcomes will be at least 6 if the outcome will be 5 or 6 in the second roll;
which is shown in figure. Similarly, if in the first roll the outcome is 2 then in the
second roll the outcome can be any one from 1 to 6; but the sum of the outcomes
will be at least 6 if the outcome will be 4 or 5 or 6 in the second roll. And this
can be done similarly for the outcome of 3 in the first roll. Now, for total number
of sample space, the die is rolled 4 times so total number of sample space is 24.
And let E be the event that the total sum of values that turn up is at least 6; then
n( E ) 2 3 4 1 10 , where 2, 3, 4 corresponds for the outcome 1, 2, 3 in the
first roll; and 1 corresponds to the outcome 6 in the first roll. So required probability will be
10 24 5 12 .
Example 6.105 [XE-2012 (1 mark)]: Suppose 50% of the population of a village like oranges, 70%
of the population like apples, and 40% like both. If a person is picked at random who likes at least one
of these fruits, what is the probability that the person likes oranges?
(a) 1 8 (b) 5 12 (c) 1 2 (d) 5 8
Solution (d): Let O , A and B are the events of liking oranges only, apple only, and both oranges
and apples, respectively, as shown in figure. So from the given data, we have O B 50% ,
B A 70% and B 40% , thus we have O 10% , B 40% , A 30% . So
the probability of a person who likes at least oranges is the sum of probability
of a person who like oranges and probability of a person who like both oranges
and apples, i.e., the required probability is
10 (10 40 30) 40 (10 40 30) 50 80 5 8 .
Example 6.106 [EC-2012, EE-2012, IN-2012 (2 marks)]: A fair coin is tossed till a head appears for
the first time. The probability that the number of required tosses is odd, is
(a) 1 3 (b) 1 2 (c) 2 3 (d) 3 4
Solution (c): We have the following sequence of events: {H , TTH , TTTTH , TTTTTTH , } . As the
probability of 1st event, i.e., appearance of head, is 1 2 ; probability of the 2nd event, i.e., appearance
of head, is 1 23 ; probability of the 3rd event, i.e., appearance of head, is 1 25 ; and so on. So the total
probability is (1 2) (1 23 ) (1 25 ) (1 2) {1 (1 22 )} (1 2) (3 4) 2 3 .
Example 6.107 [AG-2013 (2 marks)]: Box 1 contains 15 balls out of which 3 are red. Box 2 contains
12 balls out of which 4 are red. If one ball is drawn at random from each box simultaneously, the
probability of getting at least one red ball is
(a) 0.07 (b) 0.47 (c) 0.53 (d) 0.75
Solution (b): O1 , R1 are the events of drawing other colour ball and red ball from bag 1; and O2 , R2
are the events of drawing other colour ball and red ball from bag 2. As one ball is drawn at random
from each box simultaneously. Let E be the event for getting at least one red ball; and F be the
event for getting zero red ball, i.e., both balls are of other colour. So P ( E ) 1 P( F ) . As
P (O1 ) 12 15 4 5 and P (O2 ) 8 12 2 3 . So P ( F ) P (O1 ) P (O2 ) 8 15 and thus the required
probability P ( E ) 1 8 15 7 15 0.47 .
[Similar question was also asked in AG-2008 (2 marks)]
Example 6.108 [CH-2013 (1 mark)]: For two rolls of a fair die, the probability of getting a 4 in the
first roll and a number less than 4 in the second roll, up to 3 digits after the decimal point is …………
Solution: Total number of sample space when two dies are rolled is 62 36 . Let E be the event for
getting a 4 in the first roll and a number less than 4 in the second roll; so elements of E are
{(4,1), (4, 2), (4,3)} n( E ) 3 . Thus the required probability is P ( E ) 3 36 0.083 .
Example 6.109 [MN-2013 (1 mark)]: Events A and B are independent but NOT mutually
exclusive. If the probabilities P( A) and P( B) are 0.5 and 0.4 respectively, then P ( A B) is
(a) 0.6 (b) 0.7 (c) 0.8 (d) 0.9
Solution (b): As the events A and B are independent so P ( A B) P ( A) P ( B ) . Thus
P ( A B ) P( A) P( B ) P ( A B) 0.5 0.4 (0.5)(0.4) 0.7 .
Example 6.110 [BT-2014 (1 mark)]: If an unbiased coin is tossed 10 times, the probability that all
outcomes are same will be ……………. 105
10
Solution: Total number of sample space when an unbiased coin is tossed 10 times is 2 . Let E be
the event of getting same outcome; then elements of E are {HHH (10 times), TTT (10 times)}
n ( E ) 2 . Thus P ( E ) 2 210 1.95 103 195 10 5 . So answer is 195.
Example 6.111 [CE-2014 (1 mark)]: A fair (unbiased) coin was tossed four times in succession and
resulted in the following outcomes: (i) Head, (ii) Head, (iii) Head, (iv) Head. The probability of
obtaining a ‘Tail’ when the coin is tossed again is
(a) 0 (b) 1 2 (c) 4 5 (d) 1 5
Solution (b): As the event of outcome of tossing a coin is an independent event. So when the coin is
tossed again, it has no relation with previous outcomes. Thus the probability of obtaining a tail is 1 2 .
Example 6.112 [CH-2014 (2 marks)]: In rolling of two fair dice, the outcome of an experiment is
considered to be the sum of the numbers appearing on the dice. The probability is highest for the
outcome of …………….
Solution: As the total sample space when two dice are rolled is 62 36 . Let Ei be the event of
getting different sum ( i 2, 3, 4, ,12 ) of the numbers appearing on the dice. So the elements of Ei
for all the possible sum is given as: E2 {(1,1)} P ( E2 ) 1 36 , E3 {(1, 2), (2,1)} P ( E3 ) 2 36 ,
E4 {(1, 3), (2, 2), (3,1)} P ( E4 ) 3 36 , E5 {(1, 4), (2, 3), (3, 2), (4,1)} P ( E5 ) 4 36 ,
E6 {(1, 5), (2, 4), (3,3), (4, 2), (5,1)} P ( E6 ) 5 36 ,
E7 {(1, 6), (2, 5), (3, 4), (4, 3), (5, 2), (6,1)} P ( E7 ) 6 36 ,
E8 {(2, 6), (3, 5), (4, 4), (5, 3), (6, 2)} P ( E8 ) 5 36 ,
E9 {(3, 6), (4,5), (5, 4), (6, 3)} P ( E9 ) 4 36 , E10 {(4, 6), (5,5), (6, 4)} P ( E10 ) 3 36 ,
E11 {(5, 6), (6,5)} P ( E11 ) 2 36 , E12 {(6, 6)} P ( E12 ) 1 36 . So we can say that the
probability is maximum if sum of the numbers on the dice is 7. So answer is 7.
Example 6.113 [CS-2014 (2 marks)]: Four fair six-sided dice are rolled. The probability that the sum
of the results being 22 is X 1296 . The value of X is …………….
Solution: Total number of sample space when four fair six-sided dice are rolled is 6 4 . As we get the
sum of numbers appeared in the dice as 22 when we get the numbers on the dice as: (1) 6, 6, 6, 4 ,
which can be arranged in (4!) (3!) 4 ways; and (2) 6, 6, 5, 5 , which can be arranged in
(4!) (2!2!) 6 ways. So total number of ways in which we get the sum of 22 when four dice are
rolled is 4 6 10 ways. Thus the required probability is 10 6 4 10 1296 . Thus X 10 .
Example 6.114 [CS-2014 (2 marks)]: Let S be a sample space and two mutually exclusive events
A and B be such that A B S . If P () denotes the probability of the event, the maximum value of
P ( A) P( B ) is …………….
Solution: Given A B S P ( A B ) 1 P ( A) P ( B ) P ( A B ) 1 . As events A and B are
mutually exclusive so P ( A B ) 0 . Thus, P ( A) P( B ) 1 P ( B) 1 P ( A) . Let P( A) x . So
we have to find the maximum value of f ( x ) x(1 x ) x x 2 . Now f ( x) 1 2 x 0 x 1 2
and at x 1 2 , we have f ( x) 2 0 ; so x 1 2 is the point of maxima for f ( x ) . Thus
maximum value of f ( x ) is (1 2) (1 2) 2 0.25 .
Example 6.115 [CS-2014 (2 marks)]: The probability that a given positive integer lying between 1
and 100 (both inclusive) is NOT divisible by 2, 3 or 5 is …………….
Solution: Let A be the event when number [1,100] is divisible by 2; so Number of integers
[1,100] that are divisible by 2 are {2, 4, 6,8, 100} n( A) 50 P( A) 0.5 . Similarly if B and
C are the events when a numbers [1,100] is divisible by 3 and 5, respectively; then
n( B) 33 P ( B) 0.33 and n(C ) 20 P (C ) 0.2 . Now if the numbers [1,100] are divisible
by both 2 and 3, then it is also divisible by 2 3 6 and thus A B {6,12,18, 96}
n( A B ) 16 P( A B) 0.16 . Similarly, n( B C ) 6 P ( B C ) 0.06 , n( A C ) 10
P( A C ) 0.1 and n( A B C ) 3 P ( A B C ) 0.03 . So
P ( A B C ) P( A) P ( B ) P (C ) P ( A B ) P ( B C ) P (C A) P ( A B C )
P( A B C ) 0.5 0.33 0.2 0.16 0.06 0.1 0.03 0.74 , which is the probability of a
number [1,100] is divisible by 2, 3 or 5. So the required probability that a number [1,100] is not
divisible by 2, 3 or 5 is P ( Ac B c C c ) 1 P ( A B C ) 1 0.74 0.26 .
Example 6.116 [EE-2014 (2 marks)]: A fair coin is tossed n times. The probability that the
difference between the number of heads and tails is (n 3) is
(a) 2
n (b) 0 (c) n C n 3 2 n (d) 2
n 3
Solution (b): Let i and j denotes the number of heads and tails when a coin is tosses n times. As
i j n , so j n i or i n j . Now let X be the difference between the number of heads and
tails, so X i j 2i n or X i j n 2 j . Now it is given that X n 3 . So Case 1:
2i n n 3 2i 2n 3 , as the RHS is always odd and LHS is always even for any positive
integral value of n . So we never get the difference between the number of heads and tails as n 3 ;
and thus the required probability is zero. Case 2: n 2 j n 3 j 1.5 , which is not possible as j
can take only positive integral value of n ; and thus the required probability is zero.
Example 6.117 [EE-2014 (1 mark)]: Consider a dice with the property that the probability of a face
with n dots showing up is proportional to n . The probability of the face with three dots showing up is
…………….
Solution: As on a dice, the number of dots vary from 1 to 6. So n varies from 1 to 6. Since
probability of a face with n dots showing up is proportional to n , so we have P ( n) kn ; now when
a die is rolled the possible outcomes are {1, 2, 3, 4, 5, 6} , and we know that the sum of the probabilities
6
of the possible outcomes must be equal to 1. So n1 P (n) 1 k (1 2 3 4 5 6) 1
k 1 21 . Thus P (3) (1 21) 3 1 7 0.14 .
Example 6.118 [ME-2014 (1 mark)]: A box contains 25 parts of which 10 are defective. Two parts
are being drawn simultaneously in a random manner from the box. The probability of both the parts
being good is
Example 6.119 [XE-2014 (1 mark)]: Ten chocolates are distributed randomly among three children
standing in a row. The probability that the first child receives exactly three chocolates is
(a) (5 211 ) 39 (b) (5 210 ) 39 (c) 1 39 (d) 1 3
Solution (b): Let the three children be denoted by C1 , C2 , C3 ; each of the 10 chocolates can be given
to either C1 or C 2 or C3 . Hence total number of ways of distributing chocolates is 310 . Number of
10
ways to give 3 chocolates out of 10 to C1 is C3 . If C1 gets exactly 3 chocolates, remaining 7
7
chocolates should be given to C 2 or C3 ; and total number of ways of doing this is 2 . Hence,
required probability is {10 C3 27 } (310 ) (5 210 ) 39 .
Example 6.120 [PI-2016 (1 mark)]: A fair coin is tossed N times. The probability that head does
not turn up in any of the tosses is
(a) (1 2) N 1 (b) 1 (1 2) N 1 (c) (1 2) N (d) 1 (1 2) N
N
Solution (c): If a fair coin is tossed N times then total number of sample space is 2 . Now only
one time head will not turn up in any of the tosses, i.e. when all the tosses shows tails. Thus required
probability is 1 2 N . Hence option (c) is correct.
Example 6.121 [EC-2016 (1 mark)]: The probability of getting a ‘head’ in a single toss of a biased
coin is 0.3. The coin is tossed repeatedly till a ‘head’ is obtained. If the toss are independent, then the
probability of getting ‘head’ for the first time in the fifth toss is _____.
Solution: As the coin is biased so required probability of getting ‘head’ for the first time in the fifth
toss is (0.7) 4 (0.3)1 0.072 .
Example 6.123 [EE-2016 (2 marks)]: Candidates were asked to come to an interview with 3 pens
each. Black, blue, green and red were the permitted pen colours that the candidate could bring. The
probability that a candidate comes with all 3 pens having the same colour is _____.
Solution: Total number of ways in which candidates select 3 pens is same as selecting 3 pens from (3
black 3 blue 3 green 3 red pens) 12 pens is 12 C3 220 ways. Now, a candidate can select
3 pens of same colour from given 4 colours of pens, which can be done in 4 ways. So the required
probability is 4 220 0.0181 .
Example 6.124 [ME-2016 (2 marks)]: Three cards drawn from a pack of 52 cards. The probability
that they are a king, a queen, and a jack is
(a) 16 5525 (b) 64 2197 (c) 3 13 (d) 8 16575
52
Solution: Three cards from 52 cards were drawn in C3 ways.
4 4 4
A king, a queen, and a jack are drawn in C1 C1 C1 ways.
4 4 4
C1 C1 C1 4 4 4 16
So required probability is 52
.
C3 22100 5525
Example 6.125 [MN-2016 (2 marks)]: Two persons P and Q toss an unbiased coin alternately on an
understanding that whoever gets the head first wins. If P starts the game, then the probability of P
winning the game is _____.
th
Solution: Let Wi be an event of winning the game by P in i chance.
If P can wins the game in 1st chance then P(W1 ) 1 2 ;
If P loses the game in 1st chance then P wins the game only if Q loses the game in 2nd chance; thus
if P loses the game in 1st chance then P can wins the game in 3rd chance then
P (W3 ) (1 2)(1 2)(1 2) 1 23 ; Similarly
If P loses the game in 3rd chance then P can wins the game in 5th chance then
P (W5 ) (1 2)(1 2)(1 2)(1 2)(1 2) 1 2 5 ; and so on
If W be an event of winning the game by P , then required probability
1 1 1 (1 2)
P (W ) P (W1 ) P (W2 ) P (W3 ) 3 5 2
0.666 .
2 2 2 1 (1 2 )
Example 6.126 [MT-2016 (2 marks)]: A coin is tossed three times. It is known that out of the three
tosses, one is a HEAD. The probability of the other two tosses also being HEADs is _____.
3
Solution: If a coin is tossed three times then total number of possible outcomes will be 2 8 . But
out of these 8 outcomes one of the element is {TTT}; as it is given that in our sample space one out of
three tosses is HEAD, so total number of sample space is 7 (which contains at least one HEAD). Now
only one element in a sample space of 7 is {HHH}. So required probability is 1 7 0.142 .
Example 6.127 [CE-2016 (1 mark)]: X and Y are two random independent events. It is known that
P ( X ) 0.40 and P ( X Y c ) 0.70 . Which one of the following is the value of P ( X Y ) ?
(a) 0.7 (b) 0.5 (c) 0.4 (d) 0.3
Solution (a): P ( X Y ) 0.7 P ( X ) P (Y ) P ( X Y ) 0.7
As X and Y are two random independent events, so P ( X Y ) P ( X ) P (Y ) 0.40 P(Y )
Thus from (i), we have 0.4 P (Y ) 0.40 P (Y ) 0.7 0.6 P(Y ) 0.3 P (Y ) 3 6 1 2
P (Y ) 1 P (Y ) 1 2 . Now P ( X Y ) P( X ) P (Y ) P ( X Y ) ; but X and Y are two
random independent events, so P ( X Y ) P( X ) P (Y ) P ( X ) P(Y ) 0.4 0.5 0.4 0.5 0.7 .
If odds in favour of an event are a : b , then the probability of the occurrence of that event is
a ( a b ) and the probability of non-occurrence of that event is b (a b) .
If odds against an event are a : b , then the probability of the occurrence of that event is b (a b)
and the probability of non-occurrence of that event is a ( a b ) .
Example 6.128: Two dice are tossed together. Find the odds in favour of the sum of the numbers is 2.
Solution: If two dice are tossed, total number of events 6 6 36 . Favourable event is (1,1) .
Number of favourable events 1 . So the odds in favour 1 (36 1) 1 35 .
Example 6.129: A party of 23 persons take their seats at a round table. Find the odds against two persons
sitting together.
Solution: P {(21)!2!} {(22)!} 1 11 1 (1 10) . So the odd against 10 :1 .
Multiplication theorems on probability: If A and B are two events associated with a random
experiment, then P ( A B) P ( A) P( B A) , if P( A) 0 or P ( A B) P ( B ) P( A B) , if P( B) 0 .
Extension of multiplication theorem: If A1 , A2 , , An are n events related to a random
experiment, then P ( A1 A2 A3 An ) P( A1 ) P ( A2 A1 ) P A3 ( A1 A2 )
P An ( A1 A2 An 1 ) , where P Ai ( A1 A2 Ai 1 ) represents the conditional
probability of the event Ai , given that the events A1 , A2 , , Ai 1 have already happened.
Example 6.130 [CS-2000 (2 marks)]: E1 and E2 are events in a probability space satisfying the
following constraints: Pr( E1 ) Pr( E2 ) , Pr( E1 E2 ) 1 , E1 and E2 are independent. The value of
Pr( E1 ) , the probability of event E1 , is
(a) 0 (b) 1 4 (c) 1 2 (d) 1
Solution (d): Let Pr( E1 ) Pr( E2 ) p (say). As E1 and E2 are independents so
2
Pr( E1 E 2 ) Pr( E1 ) Pr( E2 ) p . Now Pr[ E1 E2 ] Pr[ E1 ] Pr[ E2 ] Pr[ E1 E2 ]
1 p p p2 p 2 2 p 1 0 p 1 .
Example 6.131 [CS-1999 (2 marks)]: Consider two events E1 and E2 such that probability of E1 ,
Pr[ E1 ] 1 2 ; probability of E2 , Pr[ E2 ] 1 3 ; probability of E1 and E2 , Pr[ E1 E2 ] 1 5 . Which
of the following statements is /are true?
(a) Pr[ E1 or E2 ] 2 3 (b) Pr[ E1 E2 ] 4 5
(c) Events E1 and E2 are not independent (d) Events E1 and E2 are independent
1
1 1 19
Solution (c): As Pr[ E1 or E2 ] Pr[ E1 E2 ] Pr[ E1 ] Pr[ E2 ] Pr[ E1 E2 ] , so
2 3 5 30
option (a) is not correct. Pr[ E1 E2 ] Pr[ E1 E2 ] Pr[ E2 ] 3 5 , so option (b) is not correct. As
Pr[ E1 E2 ] Pr[ E1 ] Pr[ E2 ] so E1 and E2 are not independent and thus option (c) is correct and
option (d) is not correct.
Example 6.132 [CS-2003 (1 mark)]: Let P ( E ) denote the probability of the event E . Given
P ( A) 1 , P ( B ) 1 2 , the values of P ( A | B) and P ( B | A) respectively are
(a) 1 4 ,1 2 (b) 1 2 ,1 4 (c) 1 2 ,1 (d) 1,1 2
Solution (d): Let A and B are exhaustive events, i.e., P( A B) 1 ; so
P ( A B ) P ( A) P ( B ) P ( A B ) 1 0.5 1 0.5 .
Thus P ( A B) P ( A B) P ( B) 0.5 0.5 1 and P ( B A) P ( A B ) P ( A) 0.5 1 0.5 .
Example 6.133 [EE-2005 (1 mark)]: If P and Q are two random events, then the following is
TRUE:
(a) Independence of P and Q implies that Probability ( P Q ) 0
(b) Probability ( P Q) Probability ( P ) + Probability (Q)
(c) If P and Q are mutually exclusive, then they must be independent
(d) Probability ( P Q) Probability ( P )
Solution (d): If P and Q are two random events then (a) Independence of P and Q means
Pr( P Q ) Pr( P ) Pr(Q) , so option (a) is not correct. As
Pr( P Q ) Pr( P ) Pr(Q) Pr( P Q) Pr( P Q) Pr( P) Pr(Q) so option (b) is not correct. As
there is no relation between mutually exclusive and independent for two random variables P and Q ,
so option (c) is not correct. As Pr( P Q ) Pr(Q P ) Pr( P ) Pr( P Q ) Pr( P) .
Example 6.134 [EE-2006 (2 marks)]: Two fair dice are rolled and the sum r of the numbers turned
up is considered
(a) Pr(r 6) 1 6 (b) Pr( r 3 is an integer) 5 6
(c) Pr( r 8 | r 4 is an integer) 5 9 (d) Pr(r 6 | r 5 is an integer) 1 18
Solution: As two fair dice are rolled so total number of sample space is 62 36 . Also r is the sum of
number appeared on 1st and 2nd dice. For r 3 , the numbers appeared on the dices is {(1, 2), (2,1)} so
P ( r 3) 2 36 ; for r 4 , the numbers appeared on the dices is {(1, 3), (2, 2), (3,1)} so
P ( r 3) 3 36 ; for r 5 , the numbers appeared on the dices is {(1, 4), (2, 3), (3, 2), (4,1)} so
P ( r 4) 4 36 ; for r 6 , the numbers appeared on the dices is {(1, 5), (2, 4), (3, 3), (4, 2), (5,1)} so
P ( r 6) 5 36 ; for r 7, the numbers appeared on the dices is
{(1, 6), (2,5), (3, 4), (4, 3), (5, 2), (6,1)} so P (r 7) 6 36 . Similarly, for P ( r 8) 5 36 ;
P (r 9) 4 36 ; P( r 10) 3 36 ; P( r 11) 2 36 ; P ( r 12) 1 36 . So
P (r 6) P( r 7) P (r 8) P (r 9) P(r 10) P(r 11) P (r 12) 21 36 7 12 . So
option (a) is not correct. Now r 3 is an integer if r 3, 6, 9,12 , so
P (r 3 is an integer) P (r 3) P (r 6) P (r 9) P(r 12) 12 36 1 3 ; thus option (b) is not
correct. Now r 4 is an integer if r 4,8,12 , so
P ( r 4 is an integer) P (r 4) P(r 8) P(r 12) 9 36 1 4 and thus
P ( r 8) 5 36 5
P ( r 8 | r 4 is an integer) . Now r 5 is an integer if r 5,10 , so
P ( r 4 is an integer) 1 4 9
P (r 5 is an integer) P(r 5) P (r 10) 7 36 and thus
P ( r 6) 5 36 5
P ( r 6 | r 5 is an integer) , so option (d) is not correct.
P ( r 5 is an integer) 7 36 7
Example 6.135 [EC-2007 (2 marks)]: An examination consists of two papers, Paper 1 and Paper 2.
The probability of failing in Paper 1 is 0.3 and that in Paper 2 is 0.2. Given that a student has failed in
Paper 2, the probability of failing in Paper 1 is 0.6. The probability of a student failing in both the
papers is:
(a) 0.5 (b) 0.18 (c) 0.12 (d) 0.06
Solution (c): Probability of failing in Paper 1 is P ( A) 0.3 ; Probability of failing in Paper 2 is
P ( B) 0.2 ; Probability of failing in Paper 1, when student has failed in Paper 2 is P ( A B ) 0.6 . So
Probability of a student failing in both the papers is P ( A B ) P( B ) P ( A B ) (0.6)(0.2) 0.12 .
Example 6.137 [CH-2009 (2 marks)]: A fair die is rolled. Let R denote the event of obtaining a
number less than or equal to 5 and S denotes the event of obtaining an odd number. Then which
ONE of the following about the probability ( P ) is TRUE?
(a) P ( R S ) 1 (b) P ( R S ) 0 (c) P ( S R ) 1 (d) P ( S R ) 0
Solution (a): Total number of sample space when a fair dice is rolled is 6. The elements of event R is
{1, 2, 3, 4,5} and the elements of event S is {1, 3, 5} ; thus the elements of ( R S ) is {1, 3, 5} . So
P ( R ) 5 6 , P( S ) 3 6 1 2 , P ( R S ) 3 6 1 2 .
P( R S ) 1 2 P( R S ) 1 2 3
Thus P ( R S ) 1 and P ( S R ) .
P(S ) 12 P( R ) 56 5
Example 6.138 [EC-2009 (2 marks)]: Consider two independent random variables X and Y with
identical distributions. The X and Y take values 0, 1 and 2 with probabilities 1 2 , 1 4 and 1 4
respectively. What is the conditional probability P ( X Y 2 | X Y 0) ?
(a) 0 (b) 1/16 (c) 1/6 (d) 1
Solution (c): We have P ( X 0) P (Y 0) 1 2 , P ( X 1) P(Y 1) 1 4 and
P ( X 2) P(Y 2) 1 2 . Now let E be the event for which X Y 2 , so the elements of E are
{( X 0, Y 2), ( X 1, Y 1), ( X 2, Y 0)} ; and F be the event for which X Y 0 , so elements
of F are {( X 0, Y 0), ( X 1, Y 1), ( X 2, Y 2)} . Thus
P ( E ) P ( X 0 and Y 2) P( X 1 and Y 1) P( X 2 and Y 0) and
P ( F ) P ( X 0 and Y 0) P( X 1 and Y 1) P ( X 2 and Y 2)
1 1 1 1 1 1 5 1 1 1 1 1 1 6
P( E ) and P ( F ) . Now event ( E F )
2 4 4 4 4 2 16 2 2 4 4 4 4 16
happens when X Y 2 and X Y 0 , which happens only when X 1 and Y 1 . As X 1 and
1 1 1
Y 1 are independent variable, so P ( E F ) P ( X 1 and Y 1) P ( X 1) P (Y 1) .
4 4 16
P ( E F ) 1 16 1
Thus the required probability, i.e., P ( X Y 2 | X Y 0) P ( E F ) .
P( F ) 6 16 6
Example 6.139 [EE-2010 (2 marks)]: A box contains 4 white balls and 3 red balls. In succession,
two balls are randomly selected and removed from the box. Given that the first removed ball is white,
the probability that the second removed ball is red is
(a) 1/3 (b) 3/7 (c) 1/2 (d) 4/7
Solution (b): Let R and W be the events of picking red and white balls, respectively. As white ball
is removed first and we have to find the probability of red ball removed when white ball was already
removed, i.e. to find P ( R W ) {P( R W )} {P (W )} . As both the events are independent so
P ( R W ) P ( R ) P (W ) (3 7) (4 7) 12 49 . Thus P ( R W ) (12 49) (4 7) 3 7 .
Example 6.140 [EC-2014 (2 marks)]: Parcels from sender S to receiver R pass sequentially
through two post-offices. Each post-office has a probability 1 5 of losing an incoming parcel,
independently of all other parcels. Given that a parcel is lost, the probability that it was lost by the
second post-office is …………….
Solution: Let A be the event for ‘parcel is lost’; B be the event for ‘parcel is lost by the 2 nd post
office’; From the given data, we have P ( B ) 1 5 , and to find P ( B A) {P( A B )} {P( A)} . Now
P ( A) P(parcel lost by 1st post office or parcel passed by 1st and lost by 2 nd post office) , so
P ( A) P(parcel lost by 1st post office) P(parcel passed by 1st and lost by 2 nd post office)
P ( A) (1 5) P (parcel passed by 1st post office) P (parcel lost by 2nd post office)
P ( A) (1 5) (4 5) (1 5) 9 25 . Now
P ( A B ) P (parcel passed by 1st and lost by 2 nd post office) P ( A B ) (4 5) (1 5) 4 25 .
Thus the required probability is P ( B A) (4 25) (9 25) 4 9 0.44 .
i 1 Ei S
n
and P ( Ei ) 0 for i 1, 2, , n , as shown in
Fig. 6.4. We can think of ( Ei ’s as the causes that lead to
the outcome of an experiment. The probabilities P( Ei ) ,
i 1, 2, , n are called prior probabilities. Suppose the Figure 6.4: Partitioning of Space S
experiment results in an outcome of event A , where
P( A) 0 . We have to find the probability that the observed event A was due to cause Ei , i.e., we
seek the conditional probability P( Ei A) . These probabilities are called posterior probabilities, given
by Baye’s rule as,
P ( Ei A) P( Ei ) P ( A Ei )
P ( Ei A) (6.14)
P( A) P( E1 ) P( A E1 ) P ( E2 ) P ( A E2 ) P ( En ) P ( A En )
where, value of P ( A) is from Eq. 6.13.
P( A Ei ) P ( Ei )
If E1 , E2 , , En form a partition of event A , then P ( A Ei ) 1 (
P ( Ei ) P ( Ei )
P( Ei )
A Ei Ei ). Hence from Eq. 6.14, we have P ( Ei A) .
P ( E1 ) P ( E2 ) P( En )
Example 6.141 [CS-2005 (2 marks)]: Box P has 2 red balls & 3 blue balls and box Q has 3 red balls
& 1 blue ball. A ball is selected as follows: (i) select a box (ii) choose a ball from the selected box
such that each ball in the box is equally likely to be chosen. The probabilities of selecting boxes P and
Q are 1 3 and 2 3 , respectively. Given that a ball selected in the above process is a red ball, the
probability that it came from the box P is:
(a) 4 19 (b) 5 19 (c) 2 9 (d) 19 30
Solution (a): Let R and B be the events for selecting red and black balls, respectively.
Example 6.143 [CS-2010 (2 marks)]: Consider a company that assembles computers. The
probability of a faulty assembly of any computer is p . The company therefore subjects each
computer to a testing process. This testing process gives the correct result for any computer with a
probability of q . What is the probability of a computer being declared faulty?
(a) pq (1 p)(1 q ) (b) (1 q ) p (c) (1 p ) q (d) pq
Solution (a): Let F & C be the event for a computer having
faulty and correct assembly, respectively; and F1 & C1 be the
event for a computer having faulty and correct testing
respectively. The diagram shows the situation of given
problem. So total probability of a computer being declared
faulty is: P (declared faulty) P( F ) P( F1 F ) P (C ) P( F1 C )
P(declared faulty) pq (1 p )(1 q)
Example 6.144 [IN-2009 (2 marks)]: A screening test is carried out to detect a certain disease. It is
found that 12% of the positive reports and 15% of the negative reports are incorrect. Assuming that
the probability of a person getting a positive report is 0.01, the probability that a person tested gets an
incorrect report is
(a) 0.0027 (b) 0.0173 (c) 0.1497 (d) 0.2100
Solution (c): Let E1 and E2 are the events that the report is positive and negative, respectively. Let
A be the event that report is incorrect. It is given that P ( E1 ) 0.01 ; also
P ( E1 ) P( E2 ) 1 P ( E2 ) 0.99 . It is also given that P ( A E1 ) 0.12 and P ( A E2 ) 0.15 . We
have to find the probability that a person tested gets an incorrect report, i.e., P( A) . So
P ( A) P( E1 ) P( A E1 ) P( E2 ) P ( A E2 ) 0.01 0.12 0.99 0.15 0.1497 .
Example 6.145 [PI-2010]: Two white and two black balls, kept in two bins, are arranged in four
ways as shown below. In each arrangement, a bin has to be chosen randomly and only one ball needs
to be picked randomly from the chosen bin. Which one of the following arrangements has the highest
probability for getting a white ball picked?
(a) (b)
(c) (d)
Solution (c): Let B1 and B2 are the events of
choosing a ball from bin 1 and 2, respectively.
Let W and B are the events of choosing a white
and black bass, respectively. Now suppose bin 1
contains w1 white balls and b1 black balls; and
bin 2 contains w2 white balls and b2 black balls.
Now, if a white ball is picked randomly from the chosen bin then, the probabilities that it is chosen
from first bin or second bin is shown in figure. Thus the total probability of choosing a white ball is:
P (W ) P ( B1 ) P (W B1 ) P ( B2 ) P (W B2 ) . Now for option (a), we have w1 1, b1 1, w2 1, b2 1 ;
thus P (W ) (1 2) (1 2) (1 2) (1 2) 1 2 . For option (b), we have w1 2, b1 0, w2 0, b2 2 ;
thus P (W ) (1 2) (2 2) (1 2) (0 2) 1 2 . For option (c), we have w1 1, b1 0, w2 1, b2 3 ;
thus P (W ) (1 2) (1 1) (1 2) (1 3) 2 3 . For option (d), we have w1 0, b1 1, w2 2, b2 1 ;
thus P (W ) (1 2) (0 1) (1 2) (2 3) 1 3 . Thus the arrangement shown in option (c) has the
highest probability for getting a white ball picked.
Example 6.146 [ME-2013, PI-2013 (2 marks)]: The probability that a student knows the correct
answer to a multiple choice question is 2/3. If the student does not know the answer, then the student
guesses the answer. The probability of the guessed answer being correct is 1/4. Given that the student
has answered the question correctly, the conditional probability that the student known the correct
answer is
2 3 5 8
(a) (b) (c) (d)
3 4 6 9
Solution (d): Let A be the events for which the student answers the question correctly; E1 be the
events for which the student knows the correct answer; and E2 be the events for which the student
2 1 1
guesses the correct answer. We have, P ( E1 ) , P( E2 ) , P ( A E2 )
; also if the student
3 3 4
knows the correct answer then he will answer the question correctly so P ( A E1 ) 1 ; we have to find
P ( E1 ) P ( A E1 ) (2 3)(1) 8
P( E1 A) , so from Eq. 6.14, P ( E1 A) .
P ( E1 ) P ( A E1 ) P ( E2 ) P ( A E2 ) (2 3)(1) (1 3)(1 4) 9
Example 6.147 [EC-2014 (1 mark)]: In a housing society, half of the families have a single child per
family, while the remaining half have two children per family. The probability that a child picked at
random, has a sibling is …………….
Solution: Let E1 and E2 be the event of single child family and two children family. Let A be the
event of picking a child. So from Bay’s theorem, we have
P ( E2 ) P ( A E2 ) (0.5)(2 x 3 x) 2
P ( E2 A) 0.667 , where x is
P( E1 ) P( A E1 ) P ( E2 ) P ( A E2 ) (0.5)( x 3x ) (0.5)(2 x 3x ) 3
the number of child from single child family; and 2x is the number of child from two children family.
Example 6.148 [ME-2014 (1 mark)]: A group consists of equal number of men and women. Of this
group 20% of the men and 50% of the women are unemployed. If a person is selected at random from
this group, the probability of the selected person being employed is …………...
Solution: Let M : Men , W : Women ,
U : Unemployed , E : Employed . The
situation is described in the diagram. So
the required probability of selected person
being employed is given as:
P ( E ) P ( M ) P( E M ) P(W ) P ( E W ) P ( E ) (1 2) (4 5) (1 2) (1 2) 13 20 0.65
Example 6.149 [CS-2016 (1 mark)]: Suppose that a shop has an equal number of LED bulbs of two
different types. The probability of an LED bulb lasting more than 100 hours given that it is of Type 1
is 0.7, and given that it is of Type 2 is 0.4. The probability that an LED bulb chosen uniformly at
random lasts more than 100 hours is _____.
Solution: Let A be the event of an LED bulb lasting more
than 100 hours and B be the event of an LED bulb lasting
not more than 100 hours. The following diagram shows the
situation of the given problem. So if an LED bulb chosen
uniformly at random lasts more than 100 hours, then its
probability is
P ( A) P{ A (Type1)}P (Type1) P{ A (Type 2)}P (Type 2)
P ( A) 0.7 0.5 0.4 0.5 0.55 .
Example 6.150 [XE-2016 (2 marks)]: A diagnostic test for a certain disease is 90% accurate. That is,
the probability of a person having (respectively, not having) the disease tested positive (respectively,
negative) is 0.9. Fifty percent of the population has the disease. What is the probability that a
randomly chosen person has the disease given that the person tested negative? _____.
Solution: Let H be an event for a person having the
disease; and NH be an event for a person not having
the disease. Also let ‘ ’ be an event for a person
tested ve ; and ‘ ' be an event for a person tested
ve . The following diagram shows the situation of the
given problem. So the probability that a randomly
chosen person has the disease given that the person
tested negative is given by:
P( H ) P ( H ) 0.1 0.5
P ( H ) P ( H ) 0.1 .
P( H ) P ( H ) P ( NH ) P( NH ) 0.1 0.5 0.9 0.5
7 6 3 7 3 6 3 7 6
P (2 B 1R) P(2 B 1R ) 0.175 0.175 0.175 0.525
10 9 8 10 9 8 10 9 8
Exercise: 6.3
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. Two fair dice are tossed. Let A be the event that the first die shows an even number and B be the
event that second die shows an odd number. The two events A and B are
(a) Mutually exclusive (b) Independent and mutually exclusive
(c) Dependent (d) None of these
2. The probabilities of a student getting I, II and III division in an examination are respectively 1 10 ,
3 5 and 1 4 . The probability that the student fail in the examination is _____.
3. If (1 3 p ) 3 , (1 p) 4 and (1 2 p) 2 are the probabilities of three mutually exclusive events,
then the set of all values of p is
(a) (1 3) p (1 2) (b) (1 3) p (1 2) (c) (1 2) p (2 3) (d) (1 2) p (2 3)
4. The probability that a leap year selected randomly will have 53 Sundays is
(a) 1 7 (b) 2 7 (c) 4 53 (d) 4 49
5. Three identical dice are rolled. The probability that same number will appear on each of them will
be
(a) 1 6 (b) 1 36 (c) 1 18 (d) 3 28
6. Three of the six vertices of a regular hexagon are chosen at random. The probability that the
triangle with these three vertices is equilateral, is equal to _____.
7. Three distinct numbers are selected from 100 natural number. The probability that all the three
numbers are divisible by 2 and 3 is
(a) 4 25 (b) 4 35 (c) 4 55 (d) 4 1155
8. Out of 21 tickets marked with numbers from 1 to 21, three are drawn at random. The chance that
the numbers on them are in A.P., is
(a) 10 133 (b) 9 133 (c) 9 1330 (d) None of these
9. There are four letters and four addressed envelopes. The chance that all letters are not dispatched
in the right envelope is
(a) 19 24 (b) 21 23 (c) 23 24 (d) 1 24
10. The letters of the word ‘ASSASSIN’ are written down at random in a row. The probability that no
two S occur together is
(a) 1 35 (b) 1 14 (c) 1 15 (d) None of these
11. A box contains 6 nails and 10 nuts. Half of the nails and half of the nuts are rusted. If one item is
chosen at random, what is the probability that it is rusted or is a nail
(a) 3 16 (b) 5 16 (c) 11 16 (d) 14 16
12. The probability that a man will be alive in 20 years is 3 5 and the probability that his wife will be
alive in 20 years is 2 3 . Then the probability that at least one will be alive in 20 years is
(a) 13 15 (b) 7 15 (c) 4 15 (d) None of these
13. Let A and B be two events such that P ( A) 0.3 and P ( A B) 0.8 . If A and B are
independent events, then P ( B )
(a) 5 6 (b) 5 7 (c) 3 5 (d) 2 5
14. A card is chosen randomly from a pack of playing cards. The probability that it is a black king or
queen of heart or jack is
(a) 1 52 (b) 6 52 (c) 7 52 (d) None of these
32. A determinant is chosen at random. The set of all determinants of order 2 with elements 0 or 1
only. The probability that value of the determinant chosen is positive, is _____.
33. If from each of the three boxes containing 3 white and 1 black, 2 white and 2 black, 1 white and 3
black balls. One ball is drawn at random. Then the probability that 2 white and 1 black ball will
be drawn is
(a) 13 32 (b) 1 4 (c) 1 32 (d) 3 16
34. The probability of India winning a test match against West Indies is 1 2 . Assuming independence
from match to match, the probability that in a 5 match series India’s second win occurs at the
third test is _____.
35. An unbiased die is tossed until a number greater than 4 appears. The probability that an even
number of tosses is needed is _____.
36. A man alternately tosses a coin and throws a dice beginning with the coin. The probability that he
gets a head in the coin before he gets a 5 or 6 in the dice is _____.
37. A box contains 100 tickets numbered 1, 2, ,100 . Two tickets are chosen at random. It is given
that the maximum number on the two chosen tickets is not more than 10. The minimum number
on them is 5 with probability
(a) 1 8 (b) 13 15 (c) 1 9 (d) None of these
38. A pair of fair dice is rolled together till a sum of either 5 or 7 is obtained. Then the probability that
5 comes before 7 is _____.
39. In an entrance test there are multiple choice questions. There are four possible answers to each
question of which one is correct. The probability that a student knows the answer to a question is
90%. If he gets the correct answer to a question, then the probability that he was guessing, is
(a) 37 40 (b) 1 37 (c) 36 37 (d) 1 9
40. A bag X contains 3 white balls and 2 black balls and another bag Y contains 2 white balls and 4
black balls. A bag and a ball out of it are picked at random. The probability that the ball is white
is _____.
41. The probability that in a year of the 22nd century chosen at random there will be 53 Sundays is
(a) 3 28 (b) 2 28 (c) 7 28 (d) 5 28
42. In a bolt factory, machines A, B and C manufacture respectively 25%, 35% and 40% of the total
bolts. Of their output 5, 4 and 2 percent respectively are defective bolts. A bolt is drawn at
random from the product. If the bolt drawn is found to be defective, the probability that it is
manufactured by the machine B is
(a) 28 69 (b) 7 69 (c) 32 69 (d) 11 69
43. An insurance company insured 2000 scooter drivers, 4000 car drivers and 6000 truck drivers. The
probability of an accident involving a scooter driver, car driver and a truck driver is 0.01, 0.03 and
0.15 respectively. One of the insured persons meets with an accident. What is the probability that
he is a scooter driver
(a) 1 52 (b) 1 62 (c) 2 51 (d) 1
44. In a test, an examinee either guesses or copies or knows the answer to a multiple choice question
with four choices. The probability that he makes a guess is 1 3 and the probability that he copies
the answer is 1 6 . The probability that his answer is correct, given that he copied it, is 1 8 . The
probability that he knew the answer to the question, given that he correctly answered it, is
(a) 24 27 (b) 24 29 (c) 24 31 (d) None of these
Example 6.151: Let X be a random variable defined as difference of the numbers that appear when a
pair of dice is rolled. What are the distinct values of the random variable?
Solution: When a pair of dice is rolled, there are 36 possible outcomes. The sample set is
S {(1,1), (1, 2), , (1, 6), (2,1), (2, 2), , (2, 6), (3,1), (3, 2), (6, 6)} . Now the random variable X is
the difference of the numbers. Its values are X {0,1, 2, 3, 4, 5} .
Probability Distribution: Each outcome i of an experiment has a probability P(i) associated with
it. Similarly, every value of random variable X xi is related to the outcome i of an experiment.
Hence, for every value of random variable xi , we have a unique real value P(i) associated. Thus,
every random variable X has probability P associated with it. This function P( X xi ) from the set
of all events of the sample space S is called a probability distribution of the random variable. The
probability distribution (or simply distribution) of a random variable X on a sample space S is set of
pairs X xi , P ( X xi ) for all xi x ( S ) , where P( X xi ) is the probability that X takes the
value xi .
Example 6.152: A random variable is number of tails when a coin is flipped thrice. Find probability
distribution of the random variable.
Solution: Sample space is HHH, THH, Value of random
X xi 0 1 2 3
HTH, HHT, TTH, THT, HTT, TTT. variable
The required probability distribution is, Probability P( X xi ) 1 8 3 8 3 8 1 8
Example 6.153: Suppose a fair dice is rolled twice. Find the possible values of random variable X
and its associated p.m.f., if X is the maximum of the two values appearing in two rolls.
Solution: Since the dice is fair, probability of any of the outcome of pair of numbers
(1,1), (1, 2), , (1, 6), (2,1), (2, 2), , (2, 6), (3,1), (3, 2), (6, 6) appearing in two rolls is 1 36 . Since X
can take values from 1 to 6. The probabilities of X taking these values can be calculated by adding
the probabilities of various outcomes that give the particular value of X . For example, X 4 , i.e.,
maximum of two values appearing on dice is four, can be obtained with outcomes
{(1, 4), (2, 4), (3, 4), (4, 4), (4, 3), (4, 2), (4,1))} i.e., 7 ways. Hence P{ X 4} 7 36 . So, the values and
probability distribution of the random variable X is given as,
X xi 1 2 3 4 5 6 P( xi )
p.m.f. ( P ( xi ) ) 1 36 3 36 5 36 7 36 9 36 11 36 1
Continuous Random Variable: Random variables could also be such that their set of possible
values is uncountable. Examples of such random variables are time between arrivals of two vehicles
at petrol pump or time taken for an angioplasty operation at a hospital or lifetime of a component.
Probability Density Function (p.d.f.): Like we have p.m.f. for discrete random variable, we
define p.d.f. for continuous random variable. Let X be a continuous random variable.
Function f ( x ) defined for all real x ( , ) is called probability density function (p.d.f.) if
for any set B of real numbers, we get probability, P{ X B} f ( x) dx . All probability
B
statements about X can be answered in terms of f ( x ) . Thus,
b
P{a X b} f ( x ) dx (6.16)
a
Note that probability of a continuous random variable at any particular value is zero, since
a
P{ X a} P{a X a} f ( x) dx 0 .
a
x
Example 6.154 [EC-2006 (1 mark)]: A probability density function is of the form p ( x) ke ,
x ( , ) . The value of k is
(a) 0.5 (b) 1 (c) 0.5 (d)
x
Solution (c): We know that for a continuous random variable p( x)dx 1 ke dx 1 ,
0
ke x dx ke x dx 1 {( k )e x }0 [{k ( )}e x }]0 1 ( k ) ( k ) 1 k 2
0
2 2
So P ( X 1) f ( x ) dx f ( x ) dx f ( x ) dx f ( x ) dx , as f ( x ) 0 for all x [0, 2] .
1 1 2 1
2 2
Thus P ( X 1) k (5 x 2 x ) dx k (5 2) x (2 3) x 3 k (17 6) (3 14) (17 6) 17 28 .
2 2
1 1
[Similar questions were also asked in TF-2007, TF-2014 (2 marks)]
Example 6.158 [EE-2013 IN-2013 (1 mark)]: A continuous random variable X has a probability
density function f ( x ) e x , 0 x , then P( X 1) is
(a) 0.368 (b) 0.5 (c) 0.632 (d) 1.0
Solution: As we have continuous probability distribution, so
P 1 x f ( x )dx e x dx e x 0 e1 1 e 0.368 .
1 1 1
Example 6.159 [CE-2013 (2 marks)]: Find the value of such that function f ( x ) is valid
( x 1)(2 x) for 1 x 2
probability density function f ( x )
0 otherwise
Solution: As for a continuous random variable, we f ( x)dx 1 . So
have
1 2 1 2
f ( x ) dx f ( x ) dx 1 ( x 1)(2 x ) dx 2 f ( x ) dx (0) dx 1 ( x 1)(2 x ) dx 2 (0)dx
2 2
f ( x ) dx (4 x x 2 2) dx 2 x 2 (1 3) x 3 2 x (1 6) . So f ( x ) dx 1 if 6 .
1 1
Example 6.160 [CE-2014 (1 mark)]: The probability density function of evaporation E on any day
1 5, 0 E 5 mm / day
during a year in a watershed is given by f ( E ) . The probability that E
0, otherwise
lies in between 2 and 4 mm/day in a day in the watershed is (in decimal) …………….
Solution: As we have continuous probability distribution, so
4 4 4
P 2 E 4 f ( E ) de (1 5) de (e 5) 2 (4 5) ( 2 5) 2 5 0.4 .
2 2
Example 6.161 [EE-2014 (2 marks)]: Let X be a random variable with probability density function
0.2, x 1
f ( x ) 0.1, 1 x 4 . The probability P (0.5 X 5) is …………….
0, otherwise
5
Solution: As we have continuous probability distribution, so P (0.5 x 5) f ( x )dx . As the
0.5
0.2, x [ 1,1]
given PDF can be written as f ( x ) 0.1, x [4, 1) (1, 4] . So
0, otherwise
5 1 4 5 1 4 5
P (0.5 x 5) f ( x) dx f ( x )dx f ( x ) dx f ( x) dx 0.2dx 0.1dx (0) dx
0.5 0.5 1 4 0.5 1 4
Example 6.162 [EE-2014 (1 mark)]: Lifetime of an electric bulb is a random variable with density
f ( x ) kx 2 , where x is measured in years. If the minimum and maximum lifetimes of bulb are 1 and
2 years respectively, then the value of k is …………….
Solution: As for a continuous random f ( x)dx 1 .
variable, we have So
1 2 1 2 2 3 2
f ( x)dx f ( x)dx 1 f ( x)dx 2 f ( x)dx (0)dx k 1 x dx 2 (0)dx k (1 3) x 1 . So
f ( x) dx 7 k 3 7 k 3 1 k 3 7 0.43 .
[Similar question was also asked in IN-2014 (1 mark)]
Example 6.163: Find the probability P x 1 2 for a continuous random variable whose p.d.f. is
kx (2 x), 0 x2
given as f ( x ) .
0, otherwise
Solution: For continuous random variable,
0 2 2 2 2
f ( x)dx 1 0dx 0 kx(2 x)dx 2 0dx 1 0 kx(2 x)dx 1 kx (k 2) x 3 1
0
12 12 2
k 3 4 . Hence, P x 1 2 f ( x ) dx (3 4) x (2 x ) dx (3 4) x 2 (1 4) x 3 5 32
0
Example 6.164 [CS-2016 (1 mark)]: A probability density function on the interval [ a,1] is given by
1 x 2 and outside this interval the value of the function is zero. The value of a is _____.
1 x 2 , a x 1
Solution: It is given that probability density function f ( x ) ;
0, otherwise
a 1 2
Also we have continuous random variable so f ( x)dx 1 (0)dx a (1 x ) dx (0) dx 1
1
1
0 1 x a 0 1 (1 a ) 1 1 a 0.5
Example 6.165 [MA-2016 (1 mark)]: Let the probability density function of a random variable X
x, 0 x 1 2
be f ( x ) c (2 x 1) 2 , 1 2 x 1 . Then, the value of c is equal to _____.
0, otherwise
Solution: We have continuous random variable, X , whose probability distribution function is given
0 12 1
by f ( x ) , then we must have f ( x )dx 1 f ( x )dx 0 f ( x ) dx f ( x ) dx f ( x )dx 1
12 1
0 12 1 2
(0) dx 0 xdx c (2 x 1) dx 0 dx 1
12 1
Example 6.166 [CE-2016 (2 marks)]: Probability density function of a random variable X is given
0.25, 1 x 5
below f ( x ) . P ( X 4) is
0, otherwise
(a) 3 4 (b) 1 2 (c) 1 4 (d) 1 8
4 1 4 1 4
Solution (a): P ( X 4) f ( x) dx f ( x)dx f ( x) dx 0dx 0.25dx
1 1
Example 6.167 [EE-2016 (2 marks)]: Let the probability density function of a random variable, X ,
be given as: f X ( x ) (3 2)e 3 xu ( x) ae 4 xu ( x) , where u ( x ) is the unit step function. Then the value
of ‘ a ’ and Prob { X 0} , respectively, are
(a) 2, 1 2 (b) 4, 1 2 (c) 2, 1 4 (d) 4, 1 4
0, x 0 0, x 0 1, x0
Solution (a): The unit step function, u ( x) u ( x ) u( x)
x0 1, x 0 0, x 0
1,
As we have continuous random variable, X , whose probability distribution function is given by
0
f X ( x ) , then we must have f ( x)dx 1 f ( x)dx 0 f ( x)dx 1
0
{(3 2)e3 x (0) ae4 x 3 x
(1)}dx {(3 2)e (1) ae (0)}dx 1 4x
0
0
0 4x
3 x e4 x 3 e 3 x a 1
a e dx (3 2) e dx 1 a 1 1 a 2.
0
4 2 3 0 4 2
0
0 0 3 x e4 x
4x 1 0 4x
Now P ( X 0) f X ( x ) dx {(3 2)e (0) ae (1)}dx a e dx 2 .
4 2
Example 6.168 [TF-2016 (2 marks)]: Let X be a continuous type random variable with probability
1 4 , 1 x 3
density function f ( x ) . When P ( X x ) 0.75 , the value of x is equal to ___.
0, otherwise
Solution: As we have rectangular continuous random variable, X , whose probability distribution
x x
function is given by f ( x ) , then we must have P ( X x) f ( x) dx 0.75 (1 4)dx 0.75
1 1
(1 4)( x ) x 1 0.75 x 1 3 x 2
The c.d.f. is defined for all values of x R . However, since the random variable takes only
isolated values, the c.d.f. is constant between two successive values of X and has steps at the
points xi 1, 2, , n . Thus, the c.d.f. for a discrete random variable is a step function.
F () 1 and F ( ) 0
Properties of a random variable can be studied only in terms of its c.d.f. We need not refer to
the underlying sample space or p.m.f., once we have c.d.f. of a random variable.
Example 6.169 [MN-2009 (2 marks)]: Cause-wise data of injuries in an underground coal mine for a
five-year period is given below:
Cause of Injury Fall of roof Fall of person Rope haulage Explosives Other causes
Number of Injuries 27 22 17 5 4
The cumulative probability of injury due to fall of roof and fall of person is
(a) 0.65 (b) 0.50 (c) 0.36 (d) 0.29
Solution (a): Let A be the event of falling of roof; and B be the event of falling of person. So
P ( A) 27 (27 22 17 5 4) 27 75 and P ( B ) 22 (27 22 17 5 4) 22 75 . So
cumulative probability of injury due to fall of roof and fall of person is P ( A) P ( B ) 49 75 0.65 .
P ( X a ) F ( a) lim F a (1 n)
n
b
P ( a X b ) F (b) F ( a) f ( x ) dx [This point was asked in CS-2005 (1 mark)]
a
a
For pure continuous random variable, P ( X a ) f ( x ) 0
a
P ( a X b) F (b) F ( a) P ( a) . In case of continuous random variable P ( a) 0 .
P (a X b) F (b) F ( a) P(b) . In case of continuous random variable P (b ) 0 .
P (a X b) F (b) F ( a) P (b) P (a ) . For continuous random variable P ( a) P(b) 0
P ( X a ) 1 P ( X a) 1 F ( a)
P ( X a ) 1 P ( X a ) P( X a ) 1 F (a ) P( a ) . In case of continuous random
variable P ( a) 0 .
Example 6.170 [MT-2007 (2 marks)]: The probability density function, p ( x ) , for a random
variable, x , is given by: p ( x ) (1 ) exp( x 2 ) . The probability that x lies between x1 0.6 and
x2 0.8 is [Use single step trapezoidal rule]
(a) 0 (b) 0.069 (c) 0.138 (d) 0.560
Solution (b): The probability that x lies between x1 0.6 and x2 0.8 for PDF
0.8 0.8
p ( x ) (1 ) exp( x 2 ) is I p ( x) dx (1 ) exp( x 2 ) . Now for single step trapezoidal
0.6 0.6
b
rule we have a f ( x ) dx h { f ( a ) f (b )} 2 , where h ba. Thus
0.8
I (1 ) exp( x 2 ) (0.8 0.6)(1 ) {exp( 0.82 ) exp(0.6 2 )} 2 0.069 .
0.6
a
Solution (a): From Eq. 6.18, F ( a ) f ( x ) dx , where F ( a) is the CDF; and f ( x ) is PDF. From
x 0dx, x [ 1,1]
0, x [ 1,1]
x x
given figure, f ( x ) ( x 1), 1 x 0 F ( x ) f ( x )dx ( x 1) dx, 1 x 0
( x 1), 0 x 1 x
( x 1)dx 0 x 1
0
0, x [ 1,1]
2
. Thus F ( x) (1 2) x x (1 2), 1 x 0 . So CDF has concave upward parabola for
(1 2) x 2 x (1 2), 0 x 1
1 x 0 ; and concave downward parabola for 0 x 1 ; which is shown in the figure in option (a).
Example 6.172 [TF-2008 (2 marks)]: The distribution function PX ( k ) of a random variable X with
parameter , satisfies the relation PX ( k 1) { ( k 1)}PX ( k ) , k 0,1, 2, 3, . If PX (0) e , the
expression obtained for PX ( k ) from above relation is
k k 1 k k 1
(a) PX ( k ) e (b) PX ( k ) e (c) PX ( k ) e (d) PX (k ) e
k! ( k 1)! k! (k 1)!
Solution (a): From the given data we have PX (0) e , and PX (1) PX (0 1) PX (0) e ;
0 1 1!
2 2 3 4
PX (2) PX (1 1) PX (1) e e ; Similarly PX (3) e ; PX (4) e ; and
11 2 1! 2! 3! 4!
so on. Thus we have PX ( k ) {( k ) ( k !)}e .
Example 6.173 [CS-2012 (1 mark)]: Consider a random variable X that takes values +1 and –1
with probability 0.5 each. The values of the cumulative distribution function F ( x ) at x 1 and 1
are
(a) 0 and 0.5 (b) 0 and 1 (c) 0.5 and 1 (d) 0.25 and 0.75
Solution (c): The cumulative distribution function F ( x ) P( X x ) . As the random variable X that
takes values +1 and –1 with probability 0.5 each. So F ( 1) P( X 1) P ( 1) 0.5 and
F (1) P ( X 1) P ( 1) P (1) 0.5 0.5 1 .
Expectation or Expected Value of Random Variable: One of the most important concepts in
probability theory is that of expectation of a random variable. Expected value, denoted by E ( X ) , is a
weighted average of the values of random variable, weight being the probability associated with it.
Expected value of random variable provides a central point for the distribution of values of random
variable. Thus, expected value is a mean or average value of the probability distribution of the random
variable and denoted as . Another way of interpretation justified by the ‘Strong Law of Large
Numbers’ is the average value of X that we would obtain if the random experiment is performed
infinite times. In other words, the average value of X is expected to approach ‘Expected Value’ as
trials increase infinitely.
Expected Value of Discrete Random Variable: If X is a discrete random variable with p.m.f.
P ( xi ) , the expectation of X , denoted by E ( X ) , is defined as,
n
E ( X ) i 1 xi P ( xi ) , where xi ( i 1, 2, ) are the values of X (6.20)
Example 6.174 [CS-2004 (2 marks)]: An examination paper has 150 multiple choice questions of
one mark each, with each question having four choices. Each incorrect answer fetches – 0.25 marks.
Suppose 1000 students choose all their answers randomly with uniform probability. The sum total of
the expected marks obtained by all these students is
(a) 0 (b) 2550 (c) 7525 (d) 9375
Solution (d): Let C be the event for a student having correct answer, so P (C ) 1 4 ; and W be the
event for a student having wrong answer, so P (W ) 3 4 . Let mc be the marks for correct answer;
and mw be the marks for wrong answer. So for a student expected marks for 1 question is
E (marks) P (C ) mc P(W ) mw (1 4) 1 (3 4)( 1 4) 1 16 . Thus expected marks for all
students for all questions is 1000 150 (1 16) 9375 .
Example 6.176 [MN-2007 (2 marks)]: The random variable X has the following probability mass
function: P (4) 1 4 , P (8) 1 4 , P (12) 1 4 , P (16) 1 4 . The expected value of X is
(a) 1 (b) 3 (c) 10 (d) 12
n
Solution (c): E ( X ) i 1 xi P ( xi ) 4 (1 4) 8 (1 4) 12 (1 4) 16 (1 4) 10 .
Example 6.177 [PI-2007 (2 marks)]: The random variable X takes on the values 1, 2 or 3 with
probabilities (2 5 P ) 5 , (1 3 P ) 5 and (1.5 2 P ) 5 , respectively. The values of P and E[ X ] are
respectively
(a) 0.05, 1.87 (b) 1.90, 5.87 (c) 0.05, 1.10 (d) 0.25, 1.40
Solution (a): As the sum of probabilities of occurrence of 1, 2 or 3 is equal to 1. So
(2 5P ) 5 (1 3P) 5 (1.5 2 P) 5 P 0.05 . Now
E ( X ) xi P ( xi ) 1 (2 5 0.05) 5 2 (1 3 0.05) 5 3 (1.5 2 0.05) 5 1.87 .
Example 6.178 [CS-2013 (1 mark)]: Consider an undirected random graph of eight vertices. The
probability that there is an edge between a pair of vertices is 1/2. What is the expected number of
unordered cycle of length three?
(a) 1/8 (b) 1 (c) 7 (d) 8
Solution (c): As we are looking at unordered cycles, so it suffices to choose any 3 out of 8 vertices in
any order; this can be done in 8 C3 ways. Now if they have a cycle among them, then to form a cycle
between 3 chosen vertices, we need to form edges between all pairs; so probability of this happening
n
is (1 2) (1 2) (1 2) 1 8 . Thus E ( X ) i 1 xi P ( xi ) 8C3 (1 2)3 7 .
Example 6.179 [EC-2014 (1 mark)]: Let X be a random variable which is uniformly chosen from
the set of positive odd numbers less than 100. The expectation, E[ X ] , is …………….
Solution: As the elements of X are X :{1,3, 5, 7, 99} n( X ) 50 . As for all the element, the
probability of choosing an element is 1 50 . So E ( X ) xi P ( xi )
1 1 1 1 1 50
E ( X ) 1 3 99 (1 3 99) {2 1 (50 1)2} 50 .
50 50 50 50 50 2
Example 6.180 [EC-2014 (2 marks)]: A fair coin is tossed repeatedly till both head and tail appear at
least once. The average number of tosses required is …………….
Solution: When a coin is tossed 1st time X xi P ( xi )
let the outcome will be Tail ( T ). Now
let X be the event of getting 1st Head ( H 1 12
H ) in xi i , i 1, 2, 3, , tosses; the TH 2 (1 2) (1 2) 1 2 2
required situation is given as in table. TTH 3 (1 2) (1 2) (1 2) 1 23
The average number of tosses after first
TTTH 4 (1 2) (1 2) (1 2) (1 2) 1 24
toss is the E ( X ) of the given situation,
i.e. … and so on … …
1 1 1 1 1 1 1
E ( X ) xi P( xi ) 1 2 2 3 3 …(i); E ( X ) 1 2 2 3 3 4 …(ii)
2 2 2 2 2 2 2
2 3
Now (i) – (ii) E ( X ) 1 (1 2) (1 2) (1 2 ) (1 2 ) (1 2) {1 (1 2)} 1 E ( X ) 2 .
So total number of average tosses is first toss plus E ( X ) 1 2 3
Example 6.181 [CS-1999 (1 mark)]: Suppose that the expectation of a random variable X is 5.
Which of the following statement is true?
(a) There is a sample point at which X has the value 5
(b) There is a sample point at which X has value greater than 5
(c) There is a sample point at which X has value greater than or equal to 5
(d) None of these
Solution (d): As Expected value, denoted by E ( X ) , is a weighted average of the values of random
variable, weight being the probability associated with it. So options (a), (b) and (c) are not correct.
Example 6.182 [TF-2010 (2 marks)]: If the density function of a random variable X is given by
x 2 , 0 x 2
f ( x) , then the mean value of X will be
0, otherwise
(a) 4/3 (b) 1/2 (c) 0 (d) 1/6
Solution (a): The mean value of X will be the expected value of X , i.e., from Eq. 6.22, we have
0 2 2
E ( X ) x f ( x ) dx x f ( x ) dx x f ( x ) dx x f ( x ) dx x f ( x )dx , as f ( x ) 0 for all
0 2 0
2
x (0, 2) . So E ( X ) ( x 2 2) dx ( x 3 6) 20 8 6 4 3 .
0
Properties of Expectation
Effect of change of Origin and Scale on E ( X ) : E[ aX b] aE ( X ) b
Proof: Let a random variable Y aX b , where a and b are constants and X be a random
variable. Then Y has the same probability distribution that of X .
n
For discrete random variable, E ( y ) i 1 yi p ( yi ) . But yi axi b for all i .
n n n
E ( aX b) i 1 (axi b) P( xi ) a i 1 xi P ( xi ) b i 1 P ( xi ) aE ( X ) b
n
( i 1 P ( xi ) 1 )
For continuous random variable, E (Y ) y f ( y ) dy . But y ax b for all real x and y
E ( aX b) ( ax b) f ( x ) dx a x f ( x ) dx b f ( x) dx aE ( X ) b
( f ( x ) dx 1 )
Expected value of Constant is the Constant itself. Thus, E (C ) C
n
For discrete random variable, E ( X ) i 1 xi P ( xi ) , but xi C for all i ,
n n n
E (C ) i 1 C P ( xi ) C i 1 P ( xi ) C ( i 1 P ( xi ) 1 )
For continuous random variable, E ( X ) x f ( x) dx , but x c
E (C ) C f ( x ) dx C f ( x ) dx C ( f ( x ) dx 1 )
Expectation of a function: Let, Y g ( X ) is a function of a random variable X , then Y is also
a random variable with the same probability distribution of X .
Example 6.183: A random variable is number of tails when a coin is flipped thrice. Find expectation
or mean of the random variable.
Solution: The required Random Variable X xi 0 1 2 3
probability distribution is
given in the table. Thus, p.m.f. P( X xi ) 18 38 38 18
4
xi P( xi ) 0 38 68 38
E( X ) i1 x P( x ) 3 2
i i
Example 6.184 [ME-2006 (2 marks)]: Consider a continuous random variable with probability
1 t for 1 t 0
density function f (t ) . The standard deviation of the random variable is:
1 t for 0 t 1
(a) 1 3 (b) 1 6 (c) 1 3 (d) 1 6
0 1
0 1 t 2 t3 t2 t3
Solution (b): E (T ) t f (t ) dt t (1 t ) dt t (1 t ) dt 0 . Also
1 0
2 3 1 2 3 0
0 1
2 2 0 2 t3 t4
1 2 t3 t 4 1
E (T ) t f (t ) dt t (1 t ) dt t (1 t ) dt . Thus
1 0
3 4 1 3 4 0 6
2
Var (T ) E (T 2 ) E (T ) 1 6 Standard deviation, Var (T ) 1 6.
Example 6.185 [EC-2009 (2 marks)]: A discrete random (a) Both the student and teacher are
variable X takes values from 1 to 5 with probabilities as right
shown in the table. (b) Both the student and teacher are
k 1 2 3 4 5 wrong
P( X k ) 0.1 0.2 0.4 0.2 0.1 (c) The student is wrong but the
A student calculates the mean of X as 3.5 and her teacher is right
teacher calculates the variance of X as 1.5. Which of the (d) The student is right but the
following statements is true? teacher is wrong
Solution (b): From Eq 6.20, the means for the given data is
5
E ( X ) i 1 xi P ( xi ) 1 0.1 2 0.2 3 0.4 4 0.2 5 0.1 3 .
5
E ( X 2 ) i 1 xi2 P ( xi ) 12 0.1 2 2 0.2 32 0.4 4 2 0.2 5 2 0.1 10.2 . So from Eq. 6.26,
2
we have Var ( X ) E ( X 2 ) E ( X ) 10.2 32 1.2 . Thus both are wrong.
Example 6.186 [XE-2010 (1 mark)]: The variance of the number of heads resulting from ten
independent tosses of a fair coin is
(a) 5 4 (b) 5 2 (c) 3 4 (d) 3 2
1, if head on k th toss 10
Solution (b): The number of head in the 10 coin tosses is X xk , xk th
.
k 1 0, if tail on k toss
10
As P ( xk ) 1 2 E ( xk ) xk P( xk ) 1 2 E ( X ) k 1 E ( xk ) 10 2 5 .
Also, E ( xk2 ) xk2 P ( xk ) 1 2 Var ( xk ) E ( xk2 ) {E ( xk )}2 (1 4) (1 2) 1 4 . So
10
Var ( X ) k 1V ( xk ) 10 4 5 2 .
Example 6.188 [ME-2014 (2 marks)]: A machine produces 0, 1 or 2 defective pieces in a day with
associated probability of 1 6 , 2 3 and 1 6 , respectively. The mean value and the variance of the
number of defective pieces produced by the machine in a day, respectively, are
(a) 1 and 1 3 (b) 1 3 and 1 (c) 1 and 4 3 (d) 1 3 and 4 3
Solution (a): From Eq 6.20, the means for the given data is
1 2 1
E ( X ) xi P ( xi ) 0 1 2 1 .
6 3 6
1 2 1 4 2 4 1
E ( X 2 ) xi2 P ( xi ) 02 12 22 . So, Var ( X ) E ( X 2 ) E ( X ) 12 .
6 3 6 6 3 3
Example 6.190 [CE-2016 (2 marks)]: If f ( x ) and g ( x ) are two probability density functions,
( x a) 1, a x 0 x a , a x 0
f ( x ) ( x a ) 1, 0 x a , g ( x) x a , 0 x a . Which one of the following
0, otherwise 0, otherwise
statements is true?
(a) Mean of f ( x ) and g ( x) are same; Variance of f ( x ) and g ( x ) are same
(b) Mean of f ( x ) and g ( x) are same; Variance of f ( x ) and g ( x ) are different
(c) Mean of f ( x ) and g ( x) are different; Variance of f ( x ) and g ( x ) are same
(d) Mean of f ( x ) and g ( x) are different; Variance of f ( x ) and g ( x ) are different
Solution (b):
a 0 x a x
For f ( x ) : E ( X ) x f ( x ) dx x (0) dx x 1 dx x 1 dx x (0) dx
a 0 a
a a
0 a
x2
0 a x2 x3 x 2 x3 x2
E ( X ) x dx x dx 0
a 0
a a 3a 2 a 3a 2 0
2 2 2
Variance of f ( x ) : Var ( X ) ( x ) f ( x) dx ( x 0) f ( x) dx x f ( x) dx
a 0 x a x
Var ( X ) x 2 (0) dx x 2 1 dx x 2 1 f ( x ) dx x 2 (0) dx
a 0 a
a a
0 a
x3
0 2 a 2 x3 2 x 4 x3 x 4 x3 a3
Var ( X ) x dx x x dx
a 0
a a 4a 3 a 4a 3 0 6
a 0 x a x
For g ( x ) : E ( X ) x g ( x )dx x (0) dx x dx x dx x (0) dx
a 0 a
a a
0 a
x
0 a x x3 x3 a3 a 3
E ( X ) x dx x dx 0
a 0
a a 3a a 3a 0 3 3
2 2 2
Variance of f ( x ) : Var ( X ) ( x ) g ( x) dx ( x 0) g ( x) dx x g ( x) dx
a 0 x a x
Var ( X ) x 2 (0)dx x 2 dx x 2 dx x 2 (0) dx
a 0 a
a a
0 a
x4 x 4 x3 a3
Var ( X ) . Hence option (b) is correct.
4a a 4a 3 0 2
Example 6.191 [MN-2014 (1 mark)]: The occurrence of head in a single toss of an unbiased coin is
given by a random variable X . The variance of X is …………….
Solution: The probability associated with 1 head in a single toss is 1 2 . So from Eq. 6.20, the mean is
E ( X ) xi P ( xi ) 1 (1 2) 1 2 . E ( X 2 ) xi2 P ( xi ) 12 (1 2) 1 2 . So from Eq.
2 2 2
6.25, we have Var ( X ) E ( X ) E ( X ) (1 2) (1 2) 1 4 0.25 .
The above characteristics are satisfied in the (i) Dice problem; (ii) Coin toss problem; (iii) Sampling
with replacement from a finite population; (iv) sampling with or without replacement from an infinite
(large) population. The probability distribution of a random variable X representing the number of
successes in a sequence of n Bernoulli trials, regardless of the order in which they occur, is
frequently of considerable interest. It is clear that X is a discrete random variable, assuming values
0,1, 2, , n . In order to determine its probability mass function, consider p X ( k ) , the probability of
having exactly k successes in n trials. This event can occur in as many ways as k letters S can be
placed in n boxes. Now, we have n choices for the position of the first S , n 1 choices for the
second S , … , and, finally, n k 1 choices for the position of the k th S . The total number of
possible arrangements is thus n( n 1) ( n k 1) . However, as no distinction is made of the S ’s
that are in the occupied positions, we must divide the number obtained above by the number of ways
in which k S ’s can be arranged in k boxes, that is, k ( k 1) 1 k ! . Hence, the number of ways in
n( n 1) ( n k 1) n!
which k successes can happen in n trials is n Ck . Now, for 0
k! k !( n k )!
0 n 0
success in n trials, the probability associated is p q ; similarly for 1 success in n trials, the
1 n 1
probability associated is p q ; and so on. Hence the probability associated for k successes in n
trials is p k q n k . So, the probability mass function is p X ( k ) n C k p k q n k . Hence, a binomial
experiment consists of n identical trials with probability of success p on each trial and probability of
failure q 1 p on each trial. Then, probability of k successes in n trials is given as,
n!
P ( x k ) n Ck p k q ( n k ) p k q ( n k ) , for values of k 0,1, 2, , n (6.27)
k !(n k )!
n
Mean for the Binomial Random Variable: E ( X ) k p ( k ) k 0 k n Ck p k q ( n k )
Example 6.192 [ME-1996 (2 marks)]: The probability of a defective piece being produced in a
manufacturing process is 0.01. The probability that out of 5 successive pieces, only one is defective, is
(a) (0.99) 4 (0.01) (b) (0.99)(0.01) 4 (c) 5 (0.99)(0.01) 4 (d) 5 (0.99) 4 (0.01)
Solution (d): The given experiment satisfies all the criteria for Binomial distribution. Let S be the
event for success, i.e. getting a defective piece, whose probability is p 0.01 ; and F be the event for
failure, i.e. getting a non-defective piece, whose probability is q 1 p 1 0.01 0.99 . Thus the
required probability of getting 1 defective out of 5 pieces 5 C1 (0.01)1 (0.99)5 1 5(0.01)1 (0.99) 4 .
Example 6.193 [CS-2004 (1 mark)]: If a fair coin is tossed four times. What is the probability that
two heads and two tails will result?
(a) 3 8 (b) 1 2 (c) 5 8 (d) 3 4
Solution: The given experiment satisfies all the criteria for Binomial distribution. So let S be the
event for success, i.e. getting Head; and F be the event for failure, i.e. getting Tail. Let p be the
probability of success, so p 1 2 ; and q be the probability of failure, so q 1 p 1 2 . Thus the
required probability of 2 successes in 4 trials is 4 C2 (1 2) 2 (1 2) 4 2 6 16 3 8 .
Example 6.194 [ME-2005 (1 marks)]: A lot has 10% defective items. Ten items are chosen
randomly from this lot. The probability that exactly 2 of the chosen items are defective is
(a) 0.0036 (b) 0.1937 (c) 0.2234 (d) 0.3874
Solution: The given experiment satisfies all the criteria for Binomial distribution. Let S be the event
for success, i.e. getting a defective item, whose probability is p 10 100 0.1 ; and F be the event
for failure, i.e. getting a non-defective item, whose probability is q 1 p 1 0.1 0.9 . Thus the
required probability of getting 2 defective out of 10 items is
10
C2 (0.1) 2 (0.9)10 2 45(0.1) 2 (0.9)8 0.1937 .
[Similar question was also asked in MN-2014 (2 marks)]
Example 6.195 [CS-2006 (2 marks)]: For each element in a set of size 2n , an unbiased coin is
tossed. The 2n coin tosses are independent. An element is chosen if the corresponding coin toss were
head. The probability that exactly n elements are chosen is:
(a) 2 n Cn 4 n (b) 2 n Cn 2 n (c) 1 2 n Cn (d) 1 2
Solution (a): The given experiment satisfies all the criteria for Binomial distribution. Actually the
question is asking to find the probability of n heads out of 2n coin tosses. So let S be the event for
success, i.e. getting Head; and F be the event for failure, i.e. getting Tail. Let p be the probability of
success, so p 1 2 ; and q be the probability of failure, so q 1 p 1 2 . Thus the required
probability of n successes in 2n trials is 2n
Cn (1 2) n (1 2) 2 n n 2n
Cn 2
2n 2n
Cn 4 .
n
Example 6.196 [CH-2010 (2 marks)]: X and Y are independent random variables. X follows a
binomial distribution with N 5 and p 1 2 . Y takes integer values 1 and 2, with equal probability.
Then the probability that X Y is
(a) 15 64 (b) 15 32 (c) 1 2 (d) 15 16
Solution: For two independent random variables X and Y , we have P ( X Y ) P ( X ) P (Y ) . As
Y takes integer values 1 and 2, with equal probability, the required probability such that X Y is
P ( X 1) P ( X 2) 5C1 (1 2)1 (1 2)5 1 5 C2 (1 2) 2 (1 2)5 2 (5 32) (10 32) 15 32
Example 6.197 [EC-2010 (2marks)]: A fair coin is tossed independently four times. The probability
of the event ‘the number of time heads shown up is more than the number of times tails shown up’ is
(a) 1/16 (b) 1/8 (c) 1/4 (d) 5/16
Solution: Given experiment satisfies all the criteria for Binomial distribution. As the number of time
heads shown up is more than the number of times tails shown up is done in two ways: (i) 4 Heads and
0 Tail; (ii) 3 Heads and 1 Tail. So let S be the event for success, i.e. getting Head, whose probability
p 1 2 ; and F be the event for failure, i.e. getting Tail, whose probability is q 1 p 1 2 . Thus
the required probability of for both cases is 4 C4 (1 2) 4 (1 2) 4 4 4C3 (1 2)3 (1 2) 4 3 5 16 .
Example 6.198 [MT-2010 (2 marks)]: The probability of obtaining ‘head’ n times, on tossing an
unbiased coin N times, is given by
(a) N Cn (1 2) N (b) n N (c) (1 2) N (d) N Pn (1 2) N
Solution (a): The given experiment satisfies all the criteria for Binomial distribution. Let S be the
event for success, i.e. getting Head, whose probability p 1 2 ; and F be the event for failure, i.e.
getting Tail, whose probability is q 1 p 1 2 . Thus the required probability of getting head n
times on tossing an unbiased coin N times is N
Cn (1 2) n (1 2) N n N Cn (1 2) N .
Example 6.199 [CE-2012 (2 marks)]: In an experiment, positive and negative values are equally
likely to occur. The probability of obtaining at most one negative value in five trials is
(a) 1/32 (b) 2/32 (c) 3/32 (d) 6/32
Solution: The given experiment has only two outcomes and both outcomes (positive and negative
values) are equally likely to occur; also both outcomes are independent. So we have a binomial
distribution. Let p and q are the event of getting positive and negative values, respectively. So
P ( p) 1 2 and P ( q) 1 2 . Now we have to find the probability of either 0 or 1 negative (i.e., at
most one negative value) value in five trials. So the required probability is
0 5 1 4
5 0 5 5 1 4 5! 1 1 5! 1 1 6
C0 {P (q )} {P ( p )} C0 {P ( q)} {P( p)} .
5!0! 2 2 4!1! 2 2 32
Example 6.200 [CE-2013 (1 mark)]: A 1 hour rainfall of 10 cm has return period of 50 year. The
probability that 1 hour of rainfall 10 cm or more will occur in each of two successive years is
(a) 0.04 (b) 0.20 (c) 0.02 (d) 0.0004
Solution (d): The given experiment satisfies all the criteria for Binomial distribution. Let S be the
event for success, i.e. having 1 hour of rainfall 10 cm, whose probability is probability of occurrence
of once in 50 years, so p 1 50 0.02 ; and F be the event for failure, i.e. not having 1 hour of
rainfall 10 cm or more, whose probability is q 1 p 1 0.02 0.98 . Thus the required probability
of having 1 hour of rainfall 10 cm or more will occur in each of two successive years is
2
C2 (0.02) 2 (0.98) 2 2 0.0004 .
[Similar questions were also asked in CE-1998 (2 marks), AG-2007 (1 mark)]
Example 6.201 [XE-2013 (1 mark)]: If the mean and variance of a binomial distribution are 6 and 2
respectively, then the probability of two failures is
(a) 4(2 3) 7 (b) 4(2 2 37 ) (c) 17(2 3)7 (d) 17(22 37 )
Solution (b): From the given data we have mean np 6 ; and variance 2 np (1 p ) 2 .
Solving these two, we get p 2 3 and n 9 . So the probability of two failures out of 9 trials, i.e.,
9 7 9 2
probability of 7 success out of 9 trials is P ( x 7) C7 (2 3) 1 (2 3) 4(2 3)7 .
Example 6.202 [EC-2014 (1 mark)]: An unbiased coin is tossed an infinite number of times. The
probability that the fourth head appears at the tenth toss is
(a) 0.067 (b) 0.073 (c) 0.082 (d) 0.091
Solution: As the given experiment satisfy all the criteria for Binomial distribution, so the
P (4th Head at the 10th toss) P(3 Heads in 9 tosses and 1 Head at 10th toss)
P(4th Head at the 10th toss) P (3 Heads in 9 tosses) P (1 Head at 10th toss)
P(4th Head at the 10th toss) 9
C3 (1 2)3 (1 2)93 (1 2) 21 256 0.082 .
Example 6.203 [ME-2014 (2 marks)]: Consider an unbiased cubic dice with opposite faces coloured
identically and each face coloured red, blue or green such that each colour appears only two times on
the dice. If the dice is thrown thrice, the probability of obtaining red colour on top face of the dice at
least twice is …………….
Solution: Let E be the event of appearance of red colour two times on the dice; so
p P( E ) 2 6 1 3 . So the probability of other colours appear two times on the dice is
q 1 p 2 3 . As the given experiment satisfy all the criteria for Binomial distribution, so the
required probability, i.e., if the dice is thrown thrice, the probability of obtaining red colour on top
face of the dice at least twice is
3 2 3 2 3 3 3 3
P ( x 2) P ( x 2) P ( x 3) C2 (1 3) (2 3) C3 (1 3) (2 3) (6 27) (1 27) 7 27 .
Example 6.204 [AG-2016 (2 marks)]: The maximum one day rainfall depth at 20 year return period
of a city is 150 mm. The probability of one day rainfall equal to or greater than 150 mm in the same
city occurring twice in 20 successive years is _____.
Solution: The given experiment satisfies all the criteria for Binomial distribution. Let S be the event
for success, i.e. having 1 day rainfall of 150 mm, whose probability is probability of occurrence of
once in 20 years, so p 1 20 0.05 ; and F be the event for failure, i.e. not having 1 day rainfall of
150 mm, whose probability is q 1 p 1 0.05 0.95 . Thus the probability of having 1 day rainfall
of 150 mm or more will occur in twice in 20 successive years is 20
C2 (0.05) 2 (0.95) 20 2 0.188 .
Example 6.205 [ME-2016 (2 marks)]: The probability that a screw manufactured by a company is
defective is 0.1. The company sells screws in packets containing 5 screws and gives a guarantee of
replacement if one or more screws in the packet are found to be defective. The probability that a
packet would have to be replaced is _____.
Solution: Let X be the number of screws that are defective. We have to find
P ( X 1) 1 P( X 0) . Now in given question we have binomial probability distribution so
P ( X 0) 5C0 (0.1) 0 (0.9)5 . Thus P ( X 1) 1 P ( X 0) 1 5C0 (0.1)0 (0.9)5 0.409
The Poisson Probability Distribution: Another discrete random variable that has numerous
practical applications is the Poisson random variable. Its probability distribution provides a good
model for data that represent the number of occurrences of a specified event in a given unit of time or
space. Here are some examples of experiments for which the random variable x can be modelled by
the Poisson random variable:
1. The number of calls received by a technical support specialist during a given period of time
2. The number of bacteria per small volume of fluid
3. The number of customer arrivals at a checkout counter during a given minute
4. The number of machine breakdowns during a given day
5. The number of traffic accidents on a section of freeway during a given time period
In each of the above examples, x represents the number of events that occur in a period of time
or space during which an average of m such events can be expected to occur. The only
assumptions needed when one uses the Poisson distribution to model experiments such as these are
that the counts or events occur randomly and independently of one another. [This point was asked
in EC-2014 (1 mark)]. Let a be the average number of times than an event occurs in a certain period
of time or space. The probability of k occurrences of this event is given as
a k ea
, for values of k 0,1, 2,
P( x k ) (6.28)
k!
Mean of the Poisson Random Variable
e a ak
e a ak a a 2a 2 3a 3
E ( x) k p (k ) k 0 k e
k 0 k! k 1 k! 1! 2! 3!
a a a2 a a
ae 1 ae e a . Hence, the mean value of X , the number of random
1! 2!
events per unit time or space, is a (the parameter of the distribution)
Variance of the Poisson Random Variable: As with the binomial distribution is convenient to
find E X ( X 1) initially, i.e.,
ak ea a k ea
E X ( X 1) k ( k 1) p( k ) k 0 k ( k 1) 0 0 k 2 k ( k 1)
k! k!
2 3 2
a a a 2 a a a 2 a a 2
E X ( X 1) e a e 1 a e e a
0! 1! 1! 2!
E ( X 2 ) E X ( X 1) E ( X ) E ( X 2 ) a 2 a
Example 6.207 [AG-2009 (1 mark)]: The probability function value [ f ( x)] at x 3 for Poisson
distribution with mean of 2 is
(a) 0.12 (b) 0.18 (c) 0.24 (d) 0.30
Solution (b): From Eq. 6.28, we have a 2 , k 3 , thus P ( x 3) 23 e 2 3! 0.18 .
Example 6.208 [PI-2010 (1 mark)]: If a random variable X satisfies the Poisson’s distribution with
a mean value of 2, then the probability that X 2 is
(a) 2e 2 (b) 1 2e 2 (c) 3e 2 (d) 1 3e 2
Solution (d): From equation 6.28, we have a 2 , so P ( X 2) P( x 2) P ( x 3) P ( x 4) .
2 2 e 2 23 e 2 2 4 e 2 2 22 23 24 2 22 23 24
P ( X 2) e e 3 1 2
2! 3! 4! 2! 3! 4! 2! 3! 4!
2 3 4
2 2 2 2
P( X 2) e 2 (3 e 2 ) , as e 1 2 ; P ( X 2) e 2 ( 3 e 2 ) 1 3e 2 .
2! 3! 4!
Example 6.209 [PI-2011 (1 mark)]: It is estimated that the average number of events during a year is
three. What is the probability of occurrence of not more than two events over a two-year duration?
Assume that the number of events follows a Poisson distribution.
(a) 0.052 (b) 0.062 (c) 0.072 (d) 0.082
Solution: If the average number of events during a year is three; then the average number of events
during two year is 2 3 6 . So from the given data and from Eq. 6.28, we have a 6 ; and we have
6 0 e 6 61 e 6 62 e 6 6
to find P ( x 2) P ( x 0) P( x 1) P ( x 2) 25e 0.062 .
0! 1! 2!
Example 6.210 [TF-2012 (2 marks)]: The number of neps in a carded web follows Poisson
distribution with a mean of 100 per m2. The probability that there is no nep in an area of 645 cm2 is
(a) e 6.45 (b) e 6.45 (c) e 645 (d) e 645
Solution (a): As mean number of neps per m2 is 100, i.e., per cm2 is 100 10 4 0.01 ; so per 645 cm2,
the mean number of neps is 645 0.01 6.45 . So the probability that there is no nep in an area of 645
cm2 is P ( x 0) 6.450 e 6.45 (0!) e 6.45 .
Example 6.211 [CS-2013 (2 marks)]: Suppose p is number of cars per minute passing through a
certain road junction between 5 PM and 6 PM, and p has a Poisson distribution with mean 3. What is
the probability of observing fewer than 3 cars during any given minute in this interval?
(a) 8 (2e3 ) (b) 9 (2e3 ) (c) 17 (2e 3 ) (d) 26 (2e3 )
Example 6.212 [CE-2014 (2 marks)]: A traffic office imposes on an average 5 number of penalties
daily on traffic violators. Assume that the number of penalties on different days is independent and
follows a Poisson distribution. The probability that there will be less than 4 penalties in a day is ……
Solution: From equation 6.28, we have a 5 . So required probability is
0 5
5 e 51 e 5 5 2 e 5 53 e 5
P ( X 4) P ( x 0) P ( x 1) P ( x 2) P ( x 3) 0.265 .
0! 1! 2! 3!
Example 6.213 [CE-2014 (2 marks)]: An observer counts 240 veh/hr at a specific highway location.
Assume that the vehicle arrival at the location is Poisson distributed, the probability of having one
vehicle arriving over a 30 second time interval is …………….
Solution: We have to find the probability of having one vehicle arriving over a 30 second time
interval. As average number of vehicle counted by the observer during 30 seconds is
1 2
(240 3600) 30 2 . So from equation 6.28, we have a 2 , P( x 1) 2 e (1!) 0.27 .
Example 6.214 [ME-2014 (2 marks)]: The number of accidents occurring in a plant in a month
follows Poisson distribution with mean as 5.2. The probability of occurrence of less than 2 accidents
in the plant during a randomly selected month is
(a) 0.029 (b) 0.034 (c) 0.039 (d) 0.044
Solution (b): From equation 6.28, we have a 5.2 . So required probability is
5.2 0 e 5.2 5.21 e 5.2
P ( X 2) P ( x 0) P ( x 1) 0.034 .
0! 1!
Example 6.215 [EC-2016 (1 mark)]: The second moment of a Poisson-distributed random variable
is 2. The mean of the random variable is _____.
Solution: For Poisson-distributed random variable, ( X ), the variance, Var ( X ) , of the distribution is
the mean, E ( X ) , of that distribution. So for Poisson-distribution, Var ( X ) E ( X ) x (say); and it is
given that E ( X 2 ) 2 . We know that Var ( X ) E ( X 2 ) {E ( X )}2 x 2 x 2 x 2 x 2 0
x 1, 2 , but x 2 , so x 1 .
1 (b a ) , a x b
f ( x) (6.31)
0, otherwise
where, a and b are the parameters of the distribution. This function is shown graphically in Fig. 6.6.
x2
Since, for any continuous distribution, P ( x1 X x2 ) f ( x) dx .
x1
x2 1 x2 x1
Hence, for rectangular distribution, P ( x1 X x2 ) dx (6.32)
x1 ba ba
It is also the area of the shaded rectangle in Fig. 6.6.
Mean of Rectangular Distribution
b x b2 2 a 2 2 a b ab
E ( X ) x f ( x ) dx dx
a ba ba 2 2
The Variance of Rectangular Distribution
2
2 2 b x b3 3 a 3 3 b 2 ab a 2
Since, E ( X ) x f ( x ) dx dx
a ba ba 3
2
2 2 2 b 2 ab a 2 2 2 2
a b 4b 4ab 4a 3a 6ab 3b
2
E( X )
3 2 12
2
2 (b a ) ba
. Hence the standard deviation for rectangular distribution is .
12 12
Example 6.216 [EE-2008 (1 mark)]: X is a uniformly distributed random variable that takes values
between 0 and 1. The value of E ( X 3 ) will be
(a) 0 (b) 1 8 (c) 1 4 (d) 1 2
1 (1 0) , 0 x 1 1, 0 x 1
Solution (c): From Eq. 6.31, we have, f ( x ) ; so from Eq.
0, otherwise 0, otherwise
0 1 1
6.32, we have E ( X 3 ) x 3 f ( x )dx x 3 (0) dx x 3 (1) dx x 3 (0)dx x3 dx 1 4 .
0 1 0
Example 6.217 [IN-2008 (2 marks)]: A random variable is uniformly distributed over the interval 2
to 10. Its variance will be
(a) 16/3 (b) 6 (c) 256/9 (d) 36
Solution (a): As the variance of uniformly distributed random variable over the interval [ a, b] is
2 (b a) 2 12 . Here a 2 , b 12 , so 2 (b a ) 2 12 (10 2) 2 12 16 3 .
Example 6.218 [ME-2009 (2 marks)]: The standard deviation of a uniformly distributed random
variable between 0 and 1 is
(a) 1 12 (b) 1 3 (c) 5 12 (d) 7 12
Solution (a): As the variance of uniformly distributed random variable over the interval [ a, b] is
2 (b a) 2 12 . Here a 0 , b 1 , so 2 (b a ) 2 12 (1 0) 2 12 1 12 . Thus standard
deviation, 1 12 .
Example 6.219 [EE-2011 (2 marks)]: A zero mean random signal is uniformly distributed between
limits a and a and its mean square value is equal to its variance. Then the r.m.s. value of the
signal is
(a) a 3 (b) a 2 (c) a 2 (d) a 3
Solution (a): As variance is mean square value of a uniformly distributed random signal between
limits a and a ; so r.m.s. value of the signal is Variance {a ( a )}2 12 a2 3 a 3.
Example 6.220 [EC-2014 (1 mark)]: Let X 1 , X 2 and X 3 be independent and identically distributed
random variables with the uniform distribution on [0,1] . The probability P{ X 1 is the largest} is ……
Solution: X 1 , X 2 and X 3 be independent and identically distributed random variables. So
P{ X 1 is the largest} 1 3 0.33 .
Example 6.221 [EC-2012, EE-2012, IN-2012 (1 mark)]: Two independent random variables X and
Y are uniformly distributed in the interval [ 1, 1] . The probability that max[ X , Y ] is less than 1/2 is
(a) 3 4 (b) 9 16 (c) 1 4 (d) 2 3
Solution (b): If max[ X , Y ] 1 2 then both X ,Y 1 2 . Thus
P max[ X , Y ] 1 2 P X 1 2 and Y 1 2 P X 1 2 P Y 1 2 , since X and Y are
1 1 2 ( 1) 3
1
independent variables. As a 1 and b 1 . So P X P 1 X and
2 2 1 ( 1) 4
1 3 1 1 1 3 3 9
similarly P Y . So P max[ X , Y ] P X P Y .
2 4 2 2 2 4 4 16
Example 6.222 [CS-2004 (2 marks)]: A point is randomly selected with uniform probability in the
X Y plane within the rectangle with corners at (0, 0) , (1, 0) , (1, 2) and (0, 2) . If p is the length of
the position vector of the point, the expected value of p 2 is
(a) 2 3 (b) 1 (c) 4 3 (d) 5 3
Solution: As the minimum value of p can be 0 (if the point chosen is (0, 0) , then length of position
vector will be 0), and maximum value can be 5 when point chosen is (1, 2) , because the point
(1, 2) is farthest from the origin. So p can vary from 0 to 5 . Now we know that
0 5 5
E ( p2 ) p 2 P( p )dp p 2 P ( p ) dp p 2 P ( p )dp p 2 P( p)dp , since from (, 0)
0 5 0
e y
2
2
dy
e z
2
2
dz
exp
2
dydz .
We shall now change the variables in this integral from y and z to the polar coordinates r and
2 r2
by letting y r cos and z r sin . Since, y 2 z 2 r 2 ; I 2 0 exp rdrd .
0
2
By putting v r 2 2 dv rdr , I 2 2 exp v dv 2 1 2 I 2 . Now
0
a b t
putting y 2 2 t ydy dt in the second integral, we get 2 e dt
2 2
b
a et a 0 a . Hence, a is the mean of Normal distribution.
2
The Variance of Normal Distribution
1 x a 2 1
E ( X ) ( x a)
2 2
exp 2
dx
b 2 2 b
xa b2 y 2 2 b2 y2 2
Putting y
2
dx bdy e y 2 bdy y ( ye ) dy
b
b 2 2
b2 y2 2 d y2 2
Using integration by parts, 2 y ( ye ) dy y ( ye ) dy dy . Since
2
dy
y2 2 y 2 2
b2 2 b2 2
2 0 e y 2 dy 2 b 2 . [ e y 2 dy 2 which is derived
2 2
earlier]. Hence the Variance of Normal distribution is b 2 and hence the standard deviation, b
x 1 x2 t 1 t2
( x) ( x) exp dx exp dt
2 2 2 2
x 1 x2 x 1 x2
( x ) ( x ) exp dx exp dx
2 2 2 2
x 1 x2 1 x2 1 x2 2
( x ) ( x ) exp dx exp dx e dx
x
2 2 2 2 2
2
e x 2 dx 2 (proved earlier). Hence ( x) ( x ) 1 ( x) 1 ( x)
lim ( x ) 1 , lim ( x ) 0
x x
A Gaussian distributed random variable with zero mean and standard deviation , then value
of its cumulative distribution function at the origin, i.e. (0) 1 2 . [This point was asked
in EC-2001 (1 mark), IN-2008 (2 marks)].
P ( Z x) 1 P ( z x ) P ( z x)
x
Proof: P ( Z x) ( x ) dx ( x ) dx ( x) dx 1 P ( Z x ) P ( Z x )
x
P ( a Z b) (b) ( a)
Probability table for Standardized Normal z 0.0 0.01 0.02 0.03 0.04
Distribution: The sample table allows us to 0.0 0.5000 5040 5080 5120 5160
read off probabilities of the form P ( Z x ) , 0.1 0.5398 5438 5478 5517 5557
i.e. the value of the integral given in Eq. 6.35 0.2 0.5793 5832 5871 5910 5948
or area under the shaded curve given in Fig. 0.3 0.6179 6217 6255 6293 6331
6.8. From the table we can identify that 0.4 0.6554 6591 6628 6664 6700
P ( z 0.32) 0.6255 , P ( z 0.5) 0.6915 . 0.5 0.6915 6950 6985 7019 7054
P ( z 0) 0.5000 , we can also note that the area under the curve of standardized normal
distribution from to is 1 and since the curve is symmetric at z 0 , i.e. the area
under left part of z 0 and the area under right part of z 0 is same, so P ( z 0) 0.5000 .
The area between z and z of the mean (i.e. between z ) constitute about
68% of the area under the Standardized Normal Distribution curve. [This point was asked in
MN-2009, AG-2012, MT-2012 (1 mark)]
The area between z 2 and z 2 of the mean (i.e. between z 2 ) constitute
about 95% of the area under the Standardized Normal Distribution curve. . [This point was
asked in MT-2011, (2 marks)]
The area between z 3 and z 3 of the mean, (i.e. between z 3 ) constitute
about 99.7% of the area under the Standardized Normal Distribution curve.
Example 6.223 [CS-2008 (2 marks)]: Let X be a random variable following normal distribution
with mean +1 and variance 4. Let Y be another normal variable with mean –1 and variance unknown.
If P ( X 1) P (Y 2) , the standard deviation of Y is
(a) 3 (b) 2 (c) 2 (d) 1
Solution: Let X & X and Y and Y are the mean and standard deviation of random variable
X and Y , respectively. So X 1 , X 4 2 , Y 1 . Firstly converting random variables X
X X X 1 Y Y Y 1
and Y to standard normal variable Z by Z and Z . As
X 2 Y Y
1 1 2 1 3
P ( X 1) P (Y 2) P Z PZ P Z 1 P Z .
2 Y Y
We know that in standard normal distribution, P ( Z 1) P ( Z 1) ; so 3 Y 1 Y 3 .
Example 6.224 [CE-2006 (2 marks)]: A class of first year B.Tech students is composed of four
batches A, B , C and D , each consisting of 30 students. It is found that the sessional marks of students
in Engineering Drawing in batch C have a mean of 6.6 and standard deviation of 2.3. The mean and
standard deviation of the marks for the entire class are 5.5 and 4.2, respectively. It is decided by the
course instructor to normalize the marks of the students of all batches to have the same mean and
standard deviation as that of the entire class. Due to this, the marks of a student in batch C are
changed from 8.5 to
(a) 6.0 (b) 7.0 (c) 8.0 (d) 9.0
Solution (d): Let C and C are the mean and standard deviation, respectively, of the students of
batch C . Let and are the mean and standard deviation, respectively, of the class. So we have
C 6.6 , C 2.3 , 5.5 , 4.2 . In order to normalize batch C to entire class, the normalized
x xC C x 5.5 8.5 6.6
score must be equated, i.e., Z Z C x 8.97 9.0 .
C 4.2 2.3
Example 6.226 [CE-2009 (2 marks)]: The standard normal probability function can be approximated
as F ( xn ) 1 1 exp( 1.7255 xn xn
0.12
) , where, xn is the standard normal deviate. If mean and
standard deviation of annual precipitation are 102 cm and 27 cm respectively, the probability that the
annual precipitation will be between 90 cm and 102 cm is
(a) 66.7% (b) 50.0% (c) 33.3% (d) 16.7%
Solution: Z will have standard normal distribution if Z ( X 102) 27 , and so we have
90 102 X 102 102 102 4 X 102
P (90 X 102) P P 0
27 27 27 9 27
P (90 X 102) P ( Z 0) P ( Z 4 9) f (0) f ( 4 9) (1 2) (1 3.005) 0.167 16.7%
Example 6.227 [CH-2010 (1 mark)]: The Maxwell-Boltzmann velocity distribution for the x
component of the velocity, at temperature T , is f (vt ) m (2 kT ) exp mvt (2kT ) . The 2
standard deviation of the distribution is
(a) 2kT m (b) kT m (c) kT m (d) kT (2m)
2
1 1 v
Solution (c): Given equation can be written as f (vt ) exp t
, so
kT m 2 2 kT m
comparing it with Eq. 6.33, we have, standard deviation, kT m .
Example 6.228 [IN-2010 (1 mark)]: The diameters of 10000 ball bearings were measured. The mean
diameter and standard deviation were found to be 10 mm and 0.05 mm respectively. Assuming
Gaussian distribution of measurements, it can be expected that the number of measurement more than
10.15 mm will be
(a) 230 (b) 115 (c) 15 (d) 2
Solution (c): Since the area between z 3 and z 3 of the mean, (i.e. between z 3 )
constitute about 99.7% of the area under the Standardized Normal Distribution curve. If 10 and
0.05 then 99.7% of the ball bearing lie between 10 0.15 , i.e. between 9.85 and 10.15. So total
number of ball bearing lie outside the range between 9.85 and 10.15 is (100 99.7) 0.3% of 10000,
which is 30. Thus that the number of measurement more than 10.15 mm will be half of 30, i.e. 15.
Example 6.229 [CE-2012 (1 mark)]: The annual precipitation data of a city is normally distributed
with mean and standard deviation as 1000 mm and 200 mm, respectively. The probability that the
annual precipitation will be more than 1200 mm is
(a) < 50% (b) 50% (c) 75% (d) 100%
Solution (a): Since the area between z and z of the mean, (i.e. between z )
constitute about 68% of the area under the Standardized Normal Distribution curve. If 1000 mm
and 200 mm then 68% of the ball bearing lie between 1000 200 mm, i.e. between 800 mm and
1200 mm. So total number of ball bearing lie outside the range between 800 mm and 1200 mm is
(100 68) 32% . So, the number of measurement more than 1200 mm will be half of 32%, i.e. 16%.
Example 6.230 [ME-2013, PI-2013 (1 mark)]: Let X be a nominal variable with mean 1 and
variance 4. The probability P ( X 0) is
(a) 0.5 (b) greater than zero and less than 0.5 (c) greater than 0.5 and less than 1 (d) 1.0
Solution (b): Let be the mean and be the standard deviation of random variable X . So
converting random variable X to standard normal variable as Z ( X ) ( X 1) 4 . So
P ( X 0) P Z (0 1) 4 P ( Z 0.25) P ( Z 0.25) P( Z 0) P (0 Z 0.25) . As
P ( Z 0) 0.5 . So P ( X 0) P ( Z 0) P (0 Z 0.25) lies between 0 and 0.25.
Example 6.231 [AG-2014 (2 marks)]: If f ( x ) is a normal distribution with mean 8 and standard
deviation 1, the value of f ( x ) for x 10 is
(a) 0.05 (b) 0.14 (c) 0.25 (d) 0.73
2
Solution (a): From Eq. 6.33, with 8 and 1 , we have f ( x ) (1 2 )e ( x 8) 2
2
f (10) (1 2 )e (10 8) 2 (1 2 )e 2 0.054 .
[Similar question was also asked in AG-2008 (1 marks)]
Example 6.232 [CE-2014 (1 mark)]: If { x} is a continuous, real valued random variable defined
over the interval ( , ) and its occurrence is defined by the density function given as:
(1 2){( x a ) b}2
f ( x ) 1 ( 2 b) e where ‘ a ’ and ‘ b ’ are the statistical attributes of the random
a 2
variable { x} . The value of the integral 1 (
2 b) e (1 2){( x a ) b} dx is
(a) 1 (b) 0.5 (c) (d) 2
Solution (b): The given function f ( x ) is Eq. 6.33 with mean a and b ; and we are asked to
find the area from to the mean a of f ( x ) . So value of the given integral is 0.5.
Example 6.233 [CH-2014 (2 marks)]: Consider the following two normal distributions
f1 ( x ) exp( x 2 ) , f 2 ( x ) {1 (2 )}exp {1 (4 )}( x 2 2 x 1) . If and denote the mean and
standard deviation, respectively, then
2 1 1 x 0 2 1
f1 ( x) exp( x ) exp 1 0 & 1 ; and
1 2 2 2 1 2 2
1 x 1 2
f 2 ( x)
1
2
exp
1 2
4
( x 2 x 1)
1
2 2
exp 2 1 & 2 2 . So
2 2
we have 1 2 and 1 2 12 22 .
Example 6.234 [EC-2014 (1 mark)]: Let X be a zero mean unit variance Gaussian random variable.
E X is equal to …………….
2
Solution: Z N (0,1) ( x ) (1 2 ) exp( x 2) , so E x x ( x ) dx
E x (1 2 )
2
x e x 2 dx (1 2 ) 0
2 2
( x)e x 2 dx xe x 2 dx .
0
In the first integral put x t dx dt , and limit change from to 0 and after applying definite
2 2
integral property we have E x 2(1 2 ) xe x 2 dx . Now let x 2 k xdx dk , and limit
0
remains same. So E x 2(1 2 ) e k dx 2 0.80 .
0
[Similar question was also asked in ME-1999 (2 marks)]
Example 6.235 [ME-2014 (1 mark)]: A nationalized bank has found that the daily balance available
in its savings accounts follows a normal distribution with a mean of Rs. 500 and a standard deviation
of Rs. 50. The percentage of savings account holders, who maintain an average daily balance more
than Rs. 500 is …………….
Solution: We have to find P ( x 500) ; converting it to standard normal distribution Z N (0,1) , we
x 500 500 500
get P P ( Z 0) it covered 50% area on RHS of standard normal
50 50
distribution. Thus percentage of savings account holders, who maintain an average daily balance
more than Rs. 500 is 50%.
Example 6.236 [AG-2016 (1 mark)]: The function f ( x ) represents a normal distribution whose
standard deviation and mean are 1 and 5, respectively. The value of f ( x ) at x 5 is
(a) 0.0 (b) 0.159 (c) 0.282 (d) 0.398
Solution (d): If X is a normally distributed with mean and standard deviation , then the
1 1 x 2
probability density function is given by f ( x ) exp . As 5 , 1 , so
2 2
1 1 5 5 2 1
f (5) exp exp 0 0.398 .
1 2 2 1 2
Example 6.237 [PI-2016 (1 mark)]: A normal random variable X has the following probability
( x 1)2
1 8
density function f X ( x ) e , x . Then 1 f X ( x ) dx
8
(a) 0 (b) 1 2 (c) 1 (1 e) (d) 1
Example 6.238 [ME-2016 (1 mark)]: The area (in percentage) under standard normal distribution
curve of random variable Z within limits from 3 to 3 is _____.
Solution: The area between z 3 and z 3 of the mean, (i.e. between z 3 ) constitute
about 99.7% of the area under the Standardized Normal Distribution curve. So answer is ‘99.7’.
Example 6.239 [TF-2016 (1 mark)]: Let X be a normally distributed random variable with mean 2
and variance 4. Then, the mean of ( X 2) 2 is equal to _____.
2
Solution: If a normally distributed variable X , having a mean and variance , is transformed
X
into a new variable Z , using the transformation Z , then that transformation is called a
standardized normal distribution with a mean of 0 and variance of 1.
In the given problem X be a normally distributed random variable with mean 2 and variance
X 2
2 4 2 ; so transforming X into Z , we get standardized normal distribution,
2
whose mean is 0. So answer is 0.
x x x 1 e x , x0
F ( x) f ( x )dx e x dx e x 1 e x F ( x )
0
0, otherwise
In order to obtain probabilities for the exponential distribution, we use the distribution function
F ( x ) 1 e x , since P ( x1 X x2 ) area to left of x1 – area to left of x2 , as shown in Fig.
6.10. Hence
x2 x1 x1 x2
P ( x1 X x2 ) F ( x2 ) F ( x1 ) (1 e ) (1 e )e e (6.37)
One of the intriguing properties of the exponential distribution is the ‘no memory’ property.
Suppose that X , the lifetime of a component which can fail suddenly, has an exponential
distribution with parameter . What is the probability that, having survived for a time x1 , it will
fail in the next small unit interval of time dx ?
P (survived time x1 and fail in dx )
P (fail in dx | survived time x1 ) . As survived time x1 is
P (survived time x1 )
x
P(fail in dx) f ( x1 )dx e 1 dx
redundant, P(fail in dx | survived time x1 ) x
dx
P( X x1 ) 1 F ( x1 ) e 1
Example 6.240 [IN-2007 (2 marks)]: Assume that the duration in minutes of a telephone
conversation follows the exponential distribution f ( x ) e x /5 5 , x 0 . The probability that the
conversation will exceed five minutes is
(a) 1 e (b) 1 (1 e) (c) 1 e 2 (d) 1 (1 e 2 )
Solution: Here 1 5 , and we have to find P (5 X ) , so from Eq. 6.37, we have
P (5 X ) F ( ) F (5) e 5 5 e 5 e 1 . We can also solve this question by integrating
given exponential distribution from 5 to , i.e., P (5 X ) (1 5)e x 5 dx ( e x 5 )5 1 e .
5
Example 6.241 [MN-2010 (1 mark)]: The variance of failure time (time to failures) of an electric
motor in shovel is 1600 hr2. If the failure time follows an exponential distribution, the expected failure
time in hr is
(a) 40 (b) 80 (c) 800 (d) 1600
Solution (a): As the variance of the given exponential distribution is 1600, i.e., 2 1 2 1600 . So
the expectation value will be E ( X ) 1 1600 40 .
Example 6.242 [MA-2014 (2 marks)]: The time to failure, in months, of light bulbs manufactured at
two plants A and B obey the exponential distribution with means 6 and 2 months respectively. Plant B
produces four times as many bulbs as plant A does. Bulbs from these plants are indistinguishable.
They are mixed and sold together. Given that a bulb purchased at random is working after 12 months,
the probability that it was manufactured at plant A is ….
Solution: Let X A and X B are the event that bulb from Plant A and B, respectively, working after 12
months; let A 1 A and B 1 B are the mean for bulb from Plant A and B, respectively; so
A 1 6 and B 1 2 . Now from Eq. 6.37, for bulb from Plant A, we have
A 12 A ( )
P (12 X A ) e e e 2 0.135 ; and for bulb from Plant B, we have
12 ()
P (12 X B ) e B e B e 6 0.0025 . As Plant B produces four times as many bulbs as
plant A does. So if the bulb is purchased is found to be working after 12 months then the required
P (12 X A ) 0.135
probability is 0.93 .
P (12 X A ) 4 P (12 X B ) 0.135 4 0.0025
Example 6.243: The mean time between successive telephone calls arriving randomly at a
switchboard is 30 seconds. What is the probability that the time between successive telephone calls
will be (a) less than 15 seconds; (b) between 15 and 30 seconds; (c) between 30 and 60 seconds; (d)
more than 60 seconds.
Solution: Since calls are arriving randomly with a mean time between calls of 30 seconds, X , the
time between calls, has an exponential distribution with a parameter such that 1 30 seconds,
therefore 1 30 calls/second. Note that is the mean number of calls per second, whereas 1 is
the mean-time, in seconds, between successive calls)
(a) P ( X 15) F (15) 1 e 15 30 0.3935
(b) P (15 X 30) F (30) F (15) e 15 30 e 30 30 0.2387
(c) P (30 X 60) F (60) F (30) e 30 30 e 60 30 0.2325
(d) P ( X 60) 1 F (60) 1 (1 e 60 30 ) 0.1353
Discrete Case: If X and Y are two discrete random variables, we define the joint probability
function of X and Y
P ( X x, Y y ) f ( x, y ) , where f ( x, y ) 0 , x y f ( x, y ) 1 (6.38)
Suppose that X can assume any one X ,Y y1 y2 yn Total
of m values x1 , x2 , , xm and Y
x1 f ( x1 , y1 ) f ( x1 , y2 ) f ( x1 , yn ) f1 ( x1 )
can assume any one of n values
y1 , y2 , , yn . Then the probability of
x2 f ( x2 , y1 ) f ( x2 , y2 ) f ( x2 , yn ) f1 ( x2 )
the event that X x j and Y yk is
xm f ( xm , y1 ) f ( xm , y2 ) f ( xm , yn )
f1 ( xm )
given by
Grand
P ( X x j , Y yk ) f ( x j , yk ) (6.39) Total f 2 ( y1 ) f 2 ( y2 ) f 2 ( yn ) Total
1
A joint probability function for X and Y can be represented by a joint probability table given in
table. The probability that X x j is obtained by adding all entries in the row corresponding to xi and
is given by
n
P ( X x j ) k 1 f ( x j , yk ) , for j 1, 2, , m (6.40)
and these are indicated by the entry totals in the extreme right hand column of the table. Similarly, the
probability that Y yk is obtainedby adding all entries in the column corresponding to yk and is
given by
m
P (Y yk ) j 1 f ( x j , yk ) , for k 1, 2, , n (6.41)
and these are indicated by the entry totals in the bottom row of the table. Because the probabilities in
Eqs. 6.40 and 6.41 are obtained from the margins of the table, we refer f1 ( x j ) and f 2 ( yk ) as the
m
marginal probability functions of X and Y , respectively. It should be noted that j 1 f1 ( x j ) 1 and
n m n
k 1 f2 ( yk ) 1 ; which can be written as j 1 k 1 f ( x j , yk ) 1 , which means that the total
probabilities of all entries is 1. The joint distribution function of X and Y is defined by
F ( x, y ) P ( X x, Y y ) u x v y f (u , v) (6.42)
Independent Random Variables: Suppose that X and Y are discrete random variables. If the
events X x and Y y are independent events for all x and y , then we say that X and Y are
independent random variables. In such a case, we have
P ( X x, Y y ) P ( X x ) P (Y y ) or equivalently f ( x, y ) f1 ( x ) f 2 ( y ) (6.47)
Conversely, if for all x and y , the joint probability function f ( x, y ) can be expressed as the product
of a function of x alone and a function of y alone (which are then the marginal probability functions
of X and Y ), X and Y are independent. If, however, f ( x, y ) cannot be so expressed, then X and
Y are dependent.
If X and Y are continuous random variables, we say that they are independent random variables if
the events X x and Y y are independent events for all x and y . In such a case we have
P ( X x, Y y ) P ( X x) P (Y y ) or equivalently F ( x, y ) F1 ( x ) F2 ( y ) (6.48)
where F1 ( x ) and F2 ( y ) are the (marginal) distribution functions of X and Y , respectively.
Conversely, X and Y are independent random variables if for all x and y , their joint distribution
function F ( x, y ) can be expressed as a product of a function of x alone and a function of y alone
(which are marginal distributions of X and Y , respectively). If, however, F ( x, y ) cannot be so
expressed, then X and Y are dependent. For continuous independent random variables, it is also true
that joint density function f ( x, y ) is the product of a function of x alone, f1 ( x ) , and a function of y
alone, f 2 ( y ) , and these are (marginal) density functions of X and Y , respectively.
f ( x, y ) P ( X x, Y y )
P{( X x ) (Y y )} (6.50)
f 2 ( x) P (Y y )
The Eqs. 6.49 and 6.50 are easily extended to the case where X , Y are continuous random variables,
i.e. the conditional density function of Y given X is given by
f ( x, y )
P{(Y y ) ( X x )} (6.51)
f1 ( x)
where, f ( x, y ) is the joint density function of X and Y , and f1 ( x ) is the marginal density function
of X . From Eq. 6.51, the probability of Y between c and d given that x X ( x dx ) as
d
P[{c Y d } {x X ( x dx)})] f ( y x )dy (6.52)
c
Example 6.245 [ME-2007 (2 marks)]: Let X and Y be two independent random variables. Which
one of the relations between expectation ( E ) , variance (Var ) and covariance (Cov) given below is
FALSE?
Example 6.247 [IN-2016 (1 mark)]: A voltage V1 is measured 100 times and its average and
standard deviation are 100 V and 1.5 V, respectively. A second voltage V2 , which is independent of
V1 , is measured 200 times and its average and standard deviation are 150 V and 2 V respectively. V3
is computed as: V3 V1 V2 . Then the standard deviation of V3 in volt is _____.
Solution: We know that if X and Y are independent random variables then
Var ( aX bY ) a 2Var ( X ) b 2Var (Y ) . We have X V1 and Y V2 , Var ( X ) V21 (1.5) 2 and
2 2 2 2
Var (Y ) V22 (2) 2 . So Var ( X Y ) 1 (1.5) 1 (2) 6.25 ; hence X Y V1 V2 6.25 2.5 .
Example 6.248 [XE-2016 (1 mark)]: A company records heights of all employees. Let X and Y
denote the errors in the average height of male and female employees, respectively. Assume that
X N (0, 4) and Y N (0, 9) and they are independent. Then the distribution of Z ( X Y ) 2 is
(a) N (0, 6.5) (b) N (0,3.25) (c) N (0, 2) (d) N (0,1)
Solution (b): If X and are two Y independent variables then (i) E ( aX bY ) aE ( X ) bE (Y ) ; and
(ii) Var ( aX bY ) a 2Var ( X ) b 2Var (Y ) . We are given E ( X ) 0 , E (Y ) 0 , Var ( X ) 4 ,
Var (Y ) 9 ; and we have to find E ( aX bY ) and Var (aX bY ) , where a b 1 2 .
So E (1 2) X (1 2)Y (1 2) E (Y ) (1 2) E (Y ) 0 ; and
Var{(1 2) X (1 2)Y } (1 2) 2 Var ( X ) (1 2) 2 Var (Y ) (1 4)(4 9) 3.25
Thus the distribution of Z ( X Y ) 2 is N (0,3.25) .
Example 6.249 [EC-2016 (2 marks)]: Two random variables X and Y are distributed according to
( x y ), 0 x 1, 0 y 1
f XY ( x, y ) . The probability P ( X Y 1) is _____.
0, otherwise
Solution: P ( X Y 1) f ( x, y ) dxdy
R
y 1
y 0
x 0
x 1 y
( x y ) dx dy
y 1 x 1 y
y 0 (1 2) x xy x 0 dy
2
P ( X Y 1)
y 1
P ( X Y 1) (1 2)(1 y ) 2 (1 y ) y dy
y 0
y 1 1
y 0 (1 2)(1 y ) dy (1 2) y (1 6) y 0 1 3 0.33
2 3
P ( X Y 1)
Exercise: 6.4
In all the following questions choose the correct option wherever option is given; fill the
calculated value, at the appropriate place, wherever it is asked to fill.
1. A coin is tossed six times. The probability of obtaining four or more heads is _____.
2. The probability that in a family of 4 children there will be at least one boy and at least one girl,
assuming that the probability of male birth is 0.5, is _____.
3. Out of 24 eggs, 6 are rotten. If 2 eggs are drawn at random, the probability distribution of the
number of rotten eggs that can be drawn is
(a) X x 0 1 2 (b) X x 0 1 2
P ( X x) 51 92 9 23 5 92 P ( X x) 9 23 51 92 5 92
(c) X x 0 1 2 (d) X x 0 1 2
P ( X x) 51 92 5 92 9 23 P ( X x) 5 92 9 23 51 92
4. 60% of the men having two wheelers put on their helmets. The probability that fewer than 4 out
of 5 will be using their helmets is _____.
5. For a continuous random variable whose probability distribution function is f ( x) (1 6) x k for
0 x 3 , then P (1 X 2) _____.
6. For a continuous random variable whose probability distribution function is
kx (2 x), 0 x2
f ( x) , P{x (1 2)}
0, otherwise
(a) 1 32 (b) 3 32 (c) 5 32 (d) 7 32
0, a 0
1 2, 0 a 1
7. The cumulative distribution function of a random variable ( X ) is F ( a ) 2 3, 1 a 2 .
11 12 , 2 a 3
1, 3 a
The value of P ( X 3) and P ( X 1) , respectively is
(a) 11 12 , 1 6 (b) 1 6 , 11 12 (c) 2 3 , 1 6 (d) 2 3 , 11 12
8. A die is tossed three times. A success is ‘getting 1 or 6’ on a toss. The mean and variance of the
number of success is, respectively,
21. A clothing store has determined that 30% of the people who enter the store will make a purchase.
Eight people enter the store during a one-hour period. The probability that exactly four people
will make a purchase is _____. The probability that at least one person will make a purchase is
_____.
(a) 0.685, 0.942 (b) 0.685, 0.136 (c) 0.136, 0.942 (d) 0.110, 0.685
22. The number of road construction projects that take place at any one time in a certain city follows a
Poisson distribution with a mean of 3. Find the probability that exactly five road construction
projects are currently taking place in this city. _____
23. The number of road construction projects that take place at any one time in a certain city follows a
Poisson distribution with a mean of 7. Find the probability that more than four road construction
projects are currently taking place in the city. _____
24. The number of traffic accidents that occur on a particular stretch of road during a month follows a
Poisson distribution with a mean of 7.6. Find the probability that less than three accidents will
occur next month on this stretch of road. _____
25. The number of traffic accidents that occur on a particular stretch of road during a month follows a
Poisson distribution with a mean of 7. Find the probability of observing exactly three accidents on
this stretch of road next month. _____
26. The number of traffic accidents that occur on a particular stretch of road during a month follows a
Poisson distribution with a mean of 6.8. Find the probability that the next two months will both
result in four accidents each occurring on this stretch of road. _____
27. A batch of 10 rocker cover gaskets contains 4 defective gaskets. If we draw samples of size 3
without replacement, from the batch of 10, find the probability that a sample contains 2 defective
gaskets. _____.
28. For Q. 27, concerning rocker cover gaskets the expectation and variance of samples containing 2
defective gaskets are, respectively,
(a) 1.2, 0.48 (b) 1.5, 0.48 (c) 1.2, 0.78 (d) 1.5, 0.78
29. In the manufacture of car tyres, a particular production process is know to yield 10 tyres with
defective walls in every batch of 100 tyres produced. From a production batch of 100 tyres, a
sample of 4 is selected for testing to destruction. The probability that the sample contains 1
defective tyre is _____. The expectation of the number of defectives in samples of size 4 is _____.
The variance of the number of defectives in samples of size 4 is _____.
(a) 0.4, 0.3, 0.33 (b) 0.5, 0.7, 0.4 (c) 0.3, 0.4, 0.33 (d) 0.4, 0.5, 0.7
30. A company buys batches of 20 components. Before a batch is accepted, 5 of the components are
selected at random from the batch and tested. The batch is rejected if more than one component in
the sample is found to be below standard. The probability that a batch which actually contains six
below-standard components is rejected is ______.
31. The continuous random variable X has pdf f ( x ) as shown in the
figure. The value of P (2.1 X 3.4) is _____. The value of E ( X )
is _____.
(a) 0.325, 3 (b) 0.5, 3 (c) 0.325, 2 (d) 0.5, 2
32. The thickness x of a protective coating applied to a conductor designed to work in corrosive
conditions follows a uniform distribution over the interval [20, 40] microns. The probability that
the coating is less than 35 microns thick, is _____.
33. If a random variable X is normally distributed with a mean of 100 and a standard deviation of
10, then P (82 X 112) is closest to
(a) 0.8849 (b) 0.0819 (c) 0.3037 (d) 0.1510
34. The amount of lemonade in a 2-litre bottle is normally distributed with a mean of 2.1 litres and a
standard deviation of 0.07 litres. The proportion of bottles that actually contain more than 2 litres
is
(a) 0.0540 (b) 0.9234 (c) 0.1790 (d) 0.0766
35. Which one of the following random variables would be normally distributed?
(a) The number of tails observed when a coin is tossed three times
(b) The number of times a die is rolled before a six is observed
a 1 (
2 b) e (1 2){( x a ) b} dx is _____.
42. If jobs arrive every 15 seconds on average, 4 per minute, what is the probability of waiting
less than or equal to 30 seconds? _____.
43. Accidents occur with a Poisson distribution at an average of 4 per week. What is the probability
that at least two weeks will elapse between accidents? _____.
44. Calls arrive at an average rate of 12 per hour. Find the probability that a call will occur in the next
5 minutes given that you have already waited 10 minutes for a call? _____
45. The average rate of job submissions in a busy computer centre is 4 per minute. If it can be
assumed that the number of submissions per minute interval is Poisson distributed, calculate the
probability that at least 15 seconds will elapse between any two jobs. _____.
46. Suppose the random variables X and Y have the join density function defined by
k (2 x y ), 2 x 6, 0 y 5
f ( x, y ) . If c 1 k , then the value of k is _____.
0, otherwise
Answer Keys
Answer Keys: Exercise: 6.1
1 2 3 4 5 6 7 8 9 10 11 12 13 14 15
96 1380 d b 72 a 30 7 18 b 12 480 28 a 1023
16 17 18 19 20 21 22 23 24 25 26 27 28 29 30
a d 375 78 896 1024 280 a 200 9 185 3 48 216 1630
31 32 33 34 35 36 37 38 39 40 41 42 43 44
40 360 b 150 64 196 485 205 51 535 20 a c a
Instructions: This test paper contains question on General Aptitude and Engineering
Mathematics. In all the following questions choose the correct option wherever option is given;
fill the calculated value, at the appropriate place, wherever it is asked to fill.
Marking Scheme
Question Number 1 to 5 6 to 10 11 to 15 16 to 20
Mark(s) 1 mark 2 marks 1 mark 2 marks
1. Astride is
(a) a verb (b) an adverb (c) a preposition (d) both (b) and (c)
2. Which one of the underlined parts of the given sentence contains an ERROR?
You would have to choose her, if you are looking for the best athlete to represent the school.
I II III IV
(a) I (b) II (c) III (d) IV
3. Choose the correct phrasal verb which is closest similar in meaning of the underlined word in the
given sentence. When Amit found out they were coming for him he lit out for the border.
(a) make out (b) make over (c) make off (d) make of
4. The first and the last sentence of the passage are in order and numbered 1 and 6. The rest of the
passage is split into 4 parts and numbered as 2, 3, 4, and 5. These 4 parts are not arranged in
proper order. Read the sentences and arrange them in a logical sequence to make a passage and
choose the correct sequence from the given options.
1. There were no finger prints anywhere.
2. First of all it was impossible even for a child to enter through the hole in the roof.
3. When the investigators tried to reconstruct the crime, they came up against facts.
4. Moreover, when the detectives tried to push a silver vase, it was found to be double the size
of the hole.
5. Again, the size of the hole was examined by the experts who said that nothing had been
passed through it.
6. These conclusions made the detectives think that it was a fake theft.
The proper sequence should be
(a) 2, 3, 4, 5 (b) 3, 2, 4, 5 (c) 5, 3, 4, 2 (d) 3, 4, 5, 2
5. Choose the correct option based on the given statements followed by their conclusions
Statements: I. All bullets are rifles II. Some rifles are swords
Conclusions: I. Some bullets are rifles II. All swords are bullets
(a) Only conclusion I follow (b) Only conclusion II follow
(c) Either conclusion I or II follows (d) Neither conclusion I nor II follows
6. Most of you probably did not see Mohan at close quarters. He had amazing qualities. One of
these qualities was that he managed to draw out the good in another person. The other person
may have had plenty of evil in him. But he somehow spotted the good and laid emphasis on the
good. The result was that the poor man had to try to be good. He could not help it. He would feel
ashamed when he did something wrong. Based on the given passage, consider the following
statements:
(i) The author assumes that most of us are not well-acquainted with Mohan’s powers
(ii) One of Mohan’s greatest qualities was that he could discover the good in another man
Which of the above statements is/are correct?
(a) Only (i) (b) Only (ii) (c) Both (i) and (ii) (d) Neither (i) nor (ii)
7. Rakhi goes 4 km West, then turns left and goes 3 km. Then, she turns right and goes 4 km and
then she turns right and goes 3 km; then she turned 495o towards her right and goes 5 km and
finally she turned 90o in clockwise direction and goes 5 km. At what distance (in km) is she from
the starting point now? _____
8. What is the last digit in 2457 754 ? _____
9. The sum of eight consecutive odd numbers is 656. The average of four consecutive even numbers
is 87. What is the sum of the smallest odd number and second largest even number? __________
10. If P and Q are running in a circular track, whose circumference is 120m, in a direction opposite
to that in which R is running, who is running at twice and thrice the speed of P and Q
respectively and on the same track. They start running from the same point. If P runs at the rate
of 3 m/sec, when (in seconds), after the start, will Q find himself equidistant and between P and
R for the first time? _____
11. If the rank of A [aij ]nn matrix is ( n 1) , then
(a) adj ( A) 0 (b) adj ( A) 0 (c) adj ( A) I n (d) None of these
12. The magnitude of directional derivative of f ( x, y , z ) xy 2 yz 3 at the point (2, 1,1) in the
direction of vector i 2 j 2 k is _____.
13. If Median k (Mode 2Mean) , then value of k is _____.
14. In Poisson distribution, the first two frequencies are 150 and 90, then the third frequency is
_____.
4 dx
15. 2 _____
1 x 9
19. A fair coin is tossed 10 times, then probability of getting at most six heads will be
(a) 49 64 (b) 51 64 (c) 53 64 (d) 55 64
d dy
20. Given the boundary value problem x ky 0 , 0 x 1 , with y (0) y (1) 0 . Then the
dx dx
solutions of the boundary value problem for k 1 (given by y1 ) and k 5 (given by y5 ) satisfy:
1 1 dy1 dy5 1 1 dy dy5
(a) 0 y1 y5 dx 0 (b) 0 dx dx
dx 0 (c) 0 y1 y5 dx 0 (d) 0 y1 y5 dx1 dx 0
dx
Instructions: This test paper contains question on General Aptitude and Engineering
Mathematics. In all the following questions choose the correct option wherever option is given;
fill the calculated value, at the appropriate place, wherever it is asked to fill.
Marking Scheme
Question Number 1 to 5 6 to 10 11 to 15 16 to 20
Mark(s) 1 mark 2 marks 1 mark 2 marks
1. Select which part of speech is the underlined word in the given sentence: Open the window and
let in the light.
(a) an adjective (b) a noun (c) a verb (d) an adverb
2. Select the appropriate option in place of underlined part of the sentence.
According to traditional Chinese medicine, people with healthy livers are said to be calm
and that they possess unerring judgment.
(a) are said to be calm and to possess (b) said to be calm and possessing
(c) have said to be calm and to possess (d) are said to be calm and possessive of
3. Which one of the following statements where underlined word is used correctly?
(a) The birth of her child was eminent, if not past due.
(b) There was no luggage standing by to indicate an eminent departure.
(c) She offered an abbreviated version of Sheela’s eminent departure.
(d) A sweep through the barn failed to reveal any eminent births.
4. Pointing to a photograph Simran says, “This man’s son’s sister is my mother-in-law.” How is the
woman’s husband related to the man in the photograph?
(a) Son-in-law (b) Brother (c) Nephew (d) Grandson
5. Arrange the given words in a meaningful sequence, and then choose the most appropriate
sequence among the options provided.
1. Wall 2. Clay 3. House 4. Room 5. Bricks
(a) 3, 4, 1, 2, 5 (b) 4, 3, 1, 5, 2 (c) 3, 4, 1, 5, 2 (d) 3, 4, 5, 2, 1
6. He dropped off to sleep. The cigarette slipped out of his mouth and burnt a great black hole in his
only shirt. The smart of the burn awoke him, and he got up, cursing under his breath, and fumbled
in the dark for a needle in order to sew up the hole. Otherwise his wife would see it in the
morning and would nag away at him for a couple of hours. But he could not find a needle. He fell
asleep again. Which one of the following statements best sums up the meaning of the above
passage?
(a) The man is extremely upset to find the shirt burnt and frantically tries to repair the damage
(b) The hole in the shirt and the wife’s anticipated nagging are minor problems, the greater one is
that the man cannot find a needle
(c) Neither the shirt hole nor the nagging nor the lack of a needle is of great consequence
(d) The man is terrified of his wife and dreads her discovering the burnt shirt
7. In an examination, Rajul got more marks than Suresh but not as many as Kiran. Kiran get more
marks than Gaurav and Manisha. Gaurav got less marks than Suresh but his marks are not the
lowest in the group. Who is the second in the descending order of marks?
(a) Kiran (b) Manisha (c) Rajul (d) Suresh
8. Suman borrowed a certain sum from Arun at a certain rate of simple interest for two years. He
lent this sum to Ravi at the same rate of interest compounded annually for the same period. At the
end of two years, he received ₹ 4200 as compound interest but paid ₹ 4000 only as simple
interest, find the rate of interest (in %). _____.
9. The odd one in the given sequence ‘1, 3, 10, 21, 64, 126, 388, 777, …’ is: _____.
10. Two loading machines each working 12 hours per day for 8 days handle 9000 tonnes of material
with an efficiency of 90%; while 3 other loading machines at an efficiency of 80% are set to
handle 12000 tonnes of material in 6 days. Find how many hours per day each should work?
_____
11. For what value of z , the function w defined by z e v ( cos u i sin u ) , w u iv , is not
analytic?
(a) z 1 (b) z 0 (c) All values of z (d) No value of z
12. Half the population of a town are consumers of an item. 100 investigators are appointed to find
the truth. Each investigator interviews 10 individuals. How many of them are expected to report
that 3 or less of the people interviewed are consumers of an item. _____
1 0 0
13. If A 0 2 0 and A3 6 A2 11A 10 I kI , then k _____.
0 0 3
14. Residue of x cos(1 x) at x 0 is
(a) 0.5 (b) 0.5 (c) 0.75 (d) 0.75
4
c (1 x ), x 1
15. A random variable X has a pdf f ( x ) . Then P x 0.5 _____.
0, otherwise
16. The equation of tangent plane to the surface 2 xz 2 3 xy 4 x 7 at the point (1, 1, 2) is
(a) ( x 1) ( y 2) 7( z 2) 0 (b) 3( x 1) ( y 1) 6( z 1) 0
(c) x 2 y 8( z 3) 0 (d) 7( x 1) 3( y 1) 8( z 2) 0
5 2
17. The Eigenvalue and its corresponding Eigen vector for the matrix A is
2 2
12 1
(a) 1 , X (b) 6 , X
6 2
12 1
(c) 6 , X (d) 1 , X
6 2
18. Using Taylor series method, the sum of values of y (0.1), y (0.2), y (0.3) , for the equation
dy dx x 2 y 2 with y (0) 1 , is _____.
19. Three identical dice are rolled. The probability that the same number will not appear on each of
them is
(a) 1 216 (b) 215 216 (c) 1 108 (d) 71 72
20. If F ( x ) f 2 {g ( x )} , g (1) 2 , g (1) 3 , f (2) 4 , f (2) 5 . Then F (1) _____.
Instructions: This test paper contains question on General Aptitude and Engineering
Mathematics. In all the following questions choose the correct option wherever option is given;
fill the calculated value, at the appropriate place, wherever it is asked to fill.
Marking Scheme
Question Number 1 to 5 6 to 10 11 to 15 16 to 20
Mark(s) 1 mark 2 marks 1 mark 2 marks
Instructions: This test paper contains question on General Aptitude and Engineering
Mathematics. In all the following questions choose the correct option wherever option is given;
fill the calculated value, at the appropriate place, wherever it is asked to fill.
Marking Scheme
Question Number 1 to 5 6 to 10 11 to 15 16 to 20
Mark(s) 1 mark 2 marks 1 mark 2 marks
1. After Rohan _____ in his notice at the office, his wife ran away with a lawyer.
(a) having handed (b) handing (c) had handed (d) was handing
2. Choose one from the given four options which the closest similar in meaning of the underlined
word in the given sentence. Given these constraints, we had no alternative but to suggest an
improvised solution.
(a) planned (b) outlined (c) prospective (d) makeshift
3. Choose the most appropriate word from the options given below to complete the following
sentence: It is impossible for a serious scholar to condone this _____ dismissal of respected
theories.
(a) astute (b) cavalier (c) sagacious (d) necessary
4. Choose the pair of give words which do not share a common similarity.
(a) Aphid : Paper (b) Moth : Wool (c) Termite : Wood (d) Locust : Plant
5. Choose the most appropriated option which is always associated with the given thing or word or
statement. Taj is in Agra; Agra is in India. Therefore, the Taj is in India.
(a) Always (b) Sometimes (c) Never (d) Often
6. One day we were becalmed among a group of small islands, most of which appeared to be
uninhabited. As soon as we were in want of fresh water, the captain sent the boat ashore to bring
off a cask or two. But we were mistaken in thinking there were no natives, for scarcely had we
drawn near to the shore 'when a band of savages rushed out of the bush and assembled on the
beach, brandishing their clubs and spears in a threatening manner. Consider the following
statements based on the above passage:
(i) The captain sent the boat to the shore to fetch some water
(ii) The savages brandished their spears in order to frighten the crew
(iii) The inhabitants of the islands were cruel people
Which of the above statements is/are correct?
(a) Only (i) (b) Only (iii) (c) Both (i) and (ii) (d) Both (ii) and (iii)
7. Professor Pant is forming five-person special group for his project. The group must contain one
team-leader, two computer engineers and two electrical engineers. P, Q and R are possible
computer engineers. R, S and T are possible team-leaders. U, V and W are possible electrical
engineers. Also, P and R prefers to work with each other in the same team. T prefers to work only
if V works. How many different possible groups, Professor Pant can make? _____.
8. Consider the following statement: (i) The ratio between Amit’s present age and his age after 8
years 4 : 5 ; (ii) The ratio between the present ages of Amit and Divya is 4 : 3 ; (iii) The ratio
between Divya’s present age and his age four years ago is 6 : 5 . To find the difference between
the present ages of Amit and Divya, which of the following option is correct?
(a) (i) or (iii) only (b) Any one of the three
(c) Any two of (i), (ii) and (iii) (d) All (i), (ii) and (iii) are required
x 2 y 2 z 2 a 2 and x z a , is
(a) a 2 2 (b) a 2 2 (c) a 2 3 (d) a 2 3
1 3 7
12. The sum and product of eigenvalues of the matrix A 2 1 1 are, respectively,
0 5 4
(a) 4 and 37 (b) 6 and 32
(c) 6 and 37 (d) 4 and 32
3 2
d y d y dy
13. The solution of the differential equation 3
2 4 4 y 0 is
dx dx dx
(a) y A cos 2 x B sin 2 x (b) y Ae x B cos 2 x C sin 2 x
(c) y Ae x B cos 2 x C sin 2 x (d) y Ae 2 x B cos 2 x C sin 2 x
14. Equation of two regression lines are y 2 x 3 , x 0.8 y 4 , then correlation coefficient is
_____.
15. A box contains 6 red, 4 white, 5 black balls. A person draws 4 balls from the box at random. Find
the probability that among the balls drawn there is at least one ball of each colour. _____
16. For what values of and , does the system of equations: 2 x 3 y 5 z 9 , 7 x 3 y 2 z 8 ,
2 x 3 y z , have infinite number of solutions?
(a) 5, 9 (b) 5, 9 (c) 5, 9 (d) 5, 9
17. If standard deviation of two variables X and Y are 5 and 6, respectively, and their covariance is
12, then coefficient of correlation between X and Y is _____.
18. If y ( 1) 5 , y (1) 8 , then using Newton’s central difference scheme, the value of ( dy dx ) x 0 is
_____.
19. Let f ( x ) be the function of period 2T 2 which is given on the interval ( 1,1) by f ( x ) 1 x 2
, the Fourier series of f ( x ) is
2 4
cos( n x ) 2
( 1) n 1
(a) f ( x )
3
2 n 1 n2
(b) f ( x )
3
n2
cos( n x)
n 1
2 4
(1) n 1 2 4
(1) n
(c) f ( x )
3
2
n2
cos(n x) (d) f ( x )
3
2
n2
cos(n x)
n 1 n 1
20. The distance between the origin and the point nearest to it on the surface z 2 1 xy is _____.
Instructions: This test paper contains question on General Aptitude and Engineering
Mathematics. In all the following questions choose the correct option wherever option is given;
fill the calculated value, at the appropriate place, wherever it is asked to fill.
Marking Scheme
Question Number 1 to 5 6 to 10 11 to 15 16 to 20
Mark(s) 1 mark 2 marks 1 mark 2 marks
who solved B was twice the number who solved C. The number of participants, who solved only
problem A, was one more than the number who solved problem A and at least one other problem.
Of all students who solved just one problem, half did not solve problem A. How many students
solved only problem B? _____
8. Three containers P , Q and R are having mixtures of water and milk in the ratio 5 :1 , 5 : 3 and
7 : 5 , respectively. If the capacities of the containers are in the ratio 5 : 4 : 5 , then find the ratio of
the milk to the water if the mixtures of all the three containers are mixed together.
(a) 115 : 53 (b) 115 : 54 (c) 54 :115 (d) 53 :115
9. Two cars one moving towards South and the other towards West, leave the same place at the same
time. The speed of one of them is greater than that of the other by 5 km/h. At the end of two
hours, they are at a distance of 50 km from each other. What is the speed (in km/h) of the bus
going slower? _____
10. The pie chart shows sources of income for an NGO. The total income is ₹ 40 crore. The bar chart
gives the expenditure incurred on various items: A – Food for poor; B – Education of illiterate; C
– Mid day meal programme; E – Eye camp expenses; F – Integrated street children programme.
Total expenditure is ₹ 39 crore.
Suppose in a particular year, grant from Central Government increases by 10%, foreign
contribution decrease by 10% and other income amounts remain same. If the expenses pattern
remains same, what is the percentage (in %) increase in ‘Food for poor’ sector? _____
11. The value of for which the system of equation x y z 6 , 4 x y z 0 ,
3 x 2 y 4 z 8 , have no solution, is _____.
12. The value of lim x ln(sin x ) is _____.
x 0
13. Possible solution for the differential equation y (t ) 2 y (t ) u (t ) with initial condition y (0 ) 0
is
(a) e 2 t u (t ) (b) (1 2)e 2 t u (t ) (c) (1 2)(1 e 2 t )u (t ) (d) (1 e 2t )u (t )
14. Let f ( z ) u ( x, y ) iv ( x , y ) , where z x iy . If u ( x, y ) x 2 y 2 2 xy . Then f ( z ) can be
expressed as
(a) 2 z 2 (1 i ) (b) 2 z 2 (1 i) (c) z 2 iz 2 (d) (1 i) z 2
15. Consider two random variables X and Y such that E ( X 2 ) E (Y 2 ) with zero mean. Two other
random variables are created: W1 X Y and W2 X Y . The coefficient of correlation
between W1 and W2 is _____.
16. The stationary point for the given function f ( x, y ) x 3 y 3 3axy are
(a) (0,1), (1, a ) (b) (0, 0), ( a, a) (c) (0, a ), (1, a ) (d) None of these
17. If the vertices of a triangle are A(0, 0) , A(1, 0) , C (0,1) . Find the value of the integral 2xdxdy
over the triangle is _____.
2z2 1
18. Evaluate dz , where C is the circle z 3 .
C ( z 1)2( z 2)