Applied - Physics - Notes 1

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Applied Physics
1. Mechanics
Classical mechanics
Forces; Newton's second law
Newton's second law relates the mass and velocity of a particle to a vector quantity
known as the force. If m is the mass of a particle and F is the vector sum of all applied
forces (i.e. the net applied force), Newton's second law states that
The quantity mv is called the momentum. Typically, the mass m is constant in time, and
Newton's law can be written in the simplified form
where is the acceleration. It is not always the case that m is independent of t.

For example, the mass of a rocket decreases as its propellant is ejected. Under such
circumstances, the above equation is incorrect and the full form of Newton's second law
must be used.
Energy
If a force F is applied to a particle that achieves a displacement Δs, the work done by the
force is the scalar quantity
If the mass of the particle is constant, and ΔWtotal is the total work done on the particle,
obtained by summing the work done by each applied force, from Newton's second law:
where Ek is called the kinetic energy. For a point particle, it is defined as
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For extended objects composed of many particles, the kinetic energy of the composite
body is the sum of the kinetic energies of the particles.
A particular class of forces, known as conservative forces, can be expressed as the

gradient of a scalar function, known as the potential energy and denoted Ep:
If all the forces acting on a particle are conservative, and Ep is the total potential energy,
obtained by summing the potential energies corresponding to each force
This result is known as conservation of energy and states that the total energy,
is constant in time. It is often useful, because many commonly encountered forces are
conservative.
SI units
SI mechanics units
edit
Name Symbol Dimensions Quantity
second (SI base unit) s s Time
metre (SI base unit) m m Distance
square metre m2 m2 Area
cubic metre m3 m3 Volume
metre per second m/s m·s−1 Velocity
metre per second squared m/s2 m·s−2 Acceleration
kilogram (SI base unit) kg kg Mass
−1
kilogram metre per second kg·m/s m·kg·s Momentum
−2
newton N m·kg·s Force
2 −1 −2
pascal Pa = N/m m ·kg·s Pressure

joule J = N·m m2·kg·s−2 Energy

newton metre N·m m2·kg·s−2 Torque
watt W = J/s m2·kg·s−3 Power
hertz Hz = 1/s s−1 Frequency
radian per second rad/s s−1 Angular velocity
radian per second squared rad/s2 s−2 Angular acceleration
kilogram square metre kg·m2 m2·kg Moment of inertia
kilogram square metre per second kg·m2/s m2·kg·s−1 Angular momentum
Orbital Energies, Kepler's Laws and Other Relationships

Kepler's Laws
Kepler's Three Laws can be used to describe the motion of the Planets:
1. Planets move in orbits that are ellipses
2. The planets move such that the line between the Sun and the Planet sweeps out
the same area in the same area in the same time no matter where in the orbit.
3. The square of the period of the orbit of a planet is proportional to the mean
distance from the Sun cubed.
The above rules were deduced empirically from the motions of the planet in the early
17th century, before Newton deduced the law of gravity and his laws of motion. When
Newton's laws are applied to the planets, Kepler's laws can be derived with certain
refinements.

Description of Orbits
Movement of a planet or satellite in an

orbit can be described with the above rules
and some simple plane geometry. The
parameters of the ellipse in terms of eccentricity will be used but it also helpful
understand the simple method of drawing an ellipse.
Energy of an Orbit
The Total energy of an object in orbit is the sum of kinetic energy (KE) and gravitational
potential energy (PE).
KE = 1/2 mv2
PE = - GMm/r
r = the distance of the orbiting body from the central object and
v = the velocity of the orbiting body
E = 1/2 mv2 - GMm/r
The semi-major axis is directly related to the

total energy of the orbit: E = - GM/2a
Incidently, the concept of circular velocity is useful in describing elliptical orbits. vc2 =
GM/a The energy also provides an expression for the velocity in orbit
E/m = GM/(2a) = 1/2 v2 - GM/r
and hence v2 = GM/a[2/r - 1/a]
but written in terms of the circular velocity
v2 = vc2[2/r - 1/a]

Keplers 2nd Law
We have already discussed Kepler's 1st Law without giving it its name. Kepler
discovered first that planets move in elliptical orbits about the Sun. The 2nd Law of
Kepler describes the relative velocity of the objects in their elliptical orbits. He
discovered that the line from the Sun to the planet swept out equal areas in equal times.
At first this does not seem very helpful but if we use a little geometery then we can use it
quantitatively. The diagram on the right illustrates the law. The area of the shaded
segment from A to B equal the area from segment C to D. Any body in the orbit around
the Sun (o) will travel from A to B in the same time that it travels from C to D. The rate
of sweeping our area by the line between Sun and orbiting object is called the Areal
Velocity, A . In one period, P, of the orbit the line sweeps out the area of the ellipse so we
can calculate this velocity from
A = (area of ellipse)/(period of ellipse) = ( a b) / P
A = (1 - e2)1/2 a2/P
Look at the diagram again; as an orbiting object goes from a to b the area swept out is
approximately the area of the triangle o-a-b. That area is equal to the isoceles triangle o-
a'-b' . The area of the later triangle can be calculate easily; that area is one half the base
(length a'-c-b') times the height (length o-c).
The height of the triangle is just the radius, r, of the orbit at point c and the base of the
triangle is the velocity perpendicular, v_, to the radius line at that same point times the
time of transit from a to b. So the rate of sweeping out area in the triangle at c is:
A = v_r/2
There are only two points in the orbit where the perpedicular velocity equals the orbit
velocity and that is a perihelion and aphelion. As a result we can relate the speed in orbit
at these two poins most easily.
rp = a (1 - e) and ra = a (1 + e)
So...
vara/2 = vprp/2 = ( a b) / P
and
va = vp( 1 + e)/(1 - e)
With a little more derivation (using Kepler's 3rd Law) we can show that
va = vc[( 1 + e)/( 1 - e )]1/2
and
va = vc[( 1 - e)/( 1 + e )]1/2
Kepler's 3rd Law - Relationship between Period and Semi-major Axis
This law was derived empirically by Kepler. He found that if the period of the planet was
given in years and the semi-major axis was given in Astronomical Units (AU) then
P2 = a3

It is easily derived for a circular orbit and the result applies to elliptical orbits when the
radius of the circle is replaced by the semi-major axis of the ellipse. The period of an
object in a circular orbit where r = a is
P = 2a/v
and hence since v = (GM/a)1/2
P = 2a3/2/(GM)1/2
This relationship is in metric units. If we transform to AU and years then we get 2/
(GM)1/2 = 1yr/ AU3/2
Newton refined Kepler's 3rd law using center of mass motion. When this is considered
then the mass, M is not just the mass of the central body (the Sun for the planets) but the
sum of the masses of both the 'central' and 'orbiting' object. In the case of the solar
system, Kepler was not too far off because the mass of the Sun is more than a thousand
times the masses of all the planets and their mass add only a small amount. So the
correct relationship is:
P = 2a3/2/(G(M+m))1/2
The idea of such a "space tether"?
It is not quite so simple, since each part of the tether will try to orbit at its own
velocity. If however the satellite could be anchored at a somewhat greater
distance, its centrifugal force will tighten the tether (and add stress to it, too).
That might work, if such a strong material existed. As far as we know, it's pure
science fiction (and was used that way in Arthur Clarke's book "Fountains of
Paradise"). No known material is strong enough.
--Navigational satellites of the GPS (Global Positioning System) move in 12 hour orbits.
What is their distance?
--The calculation is very similar to the preceding one.
We have 43200 sec = 5063.5 R*SQRT(R) = 5063.5 R3/2
R3/2 = 8.53164 R = (8.53164)2/3 = 4.175 RE
By such reasoning, Newton came to the conclusion that any two objects in the Universe
exert gravitational attraction on each other, with the force having a universal form:

The constant of proportionality G is known as the universal gravitational constant. It is

termed a "universal constant" because it is thought to be the same at all places and all
times, and thus universally characterizes the intrinsic strength of the gravitational force.
The Center of Mass for a Binary System
If you think about it a moment, it may seem a little strange that in Kepler's Laws the Sun
is fixed at a point in space and the planet revolves around it. Why is the Sun privileged?
Kepler had rather mystical ideas about the Sun, endowing it with almost god-like
qualities that justified its special place. However Newton, largely as a corollary of his
3rd Law, demonstrated that the situation actually was more symmetrical than Kepler
imagined and that the Sun does not occupy a privileged postion; in the process he
modified Kepler's 3rd Law.
Consider the diagram shown to the right. We may define a point called the center of
mass between two objects through the equations
where R is the total separation between the centers of the two objects. The center of
mass is familiar to anyone who has ever played on a see-saw. The fulcrum point at which
the see-saw will exactly balance two people sitting on either end is the center of mass for
the two persons sitting on the see-saw.
Here is a Center of Mass Calculator that will help you make and visualize calculations
concerning the center of mass. (Caution: this applet is written under Java 1.1, which is
only supported by the most recent browsers. It should work on Windows systems under
Netscape 4.06 or the most recent version of Internet Explorer 4.0, but may not yet work
on Mac or Unix systems or earlier Windows browsers.)
Newton's Modification of Kepler's Third Law

Because for every action there is an equal and opposite reaction, Newton realized that in
the planet-Sun system the planet does not orbit around a stationary Sun. Instead, Newton
proposed that both the planet and the Sun orbited around the common center of mass for
the planet-Sun system. He then modified Kepler's 3rd Law to read,
where P is the planetary orbital period and the other quantities have the meanings
described above, with the Sun as one mass and the planet as the other mass. (As in the
earlier discussion of Kepler's 3rd Law, this form of the equation assumes that masses are
measured in solar masses, times in Earth years, and distances in astronomical units.)
Notice the symmetry of this equation: since the masses are added on the left side and the

distances are added on the right side, it doesn't matter whether the Sun is labeled with 1
and the planet with 2, or vice-versa. One obtains the same result in either case.
Now notice what happens in Newton's new equation if one of the masses (either 1 or 2;
remember the symmetry) is very large compared with the other. In particular, suppose
the Sun is labeled as mass 1, and its mass is much larger than the mass for any of the
planets. Then the sum of the two masses is always approximately equal to the mass of
the Sun, and if we take ratios of Kepler's 3rd Law for two different planets the masses
cancel from the ratio and we are left with the original form of Kepler's 3rd Law:
Thus Kepler's 3rd Law is approximately valid because the Sun is much more massive
than any of the planets and therefore Newton's correction is small. The data Kepler had
access to were not good enough to show this small effect. However, detailed
observations made after Kepler show that Newton's modified form of Kepler's 3rd Law
is in better accord with the data than Kepler's original form.
Two Limiting Cases
We can gain further insight by considering the position of the center of mass in two
limits. First consider the example just addressed, where one mass is much larger than the
other. Then, we see that the center of mass for the system essentially concides with the
center of the massive object:
This is the situation in the Solar System: the Sun is so massive compared with any of the
planets that the center of mass for a Sun-planet pair is always very near the center of the
Sun. Thus, for all practical purposes the Sun IS almost (but not quite) motionless at the
center of mass for the system, as Kepler originally thought.
However, now consider the other limiting case where the two masses are equal to each
other. Then it is easy to see that the center of mass lies equidistant from the two masses
and if they are gravitationally bound to each other, each mass orbits the common center
of mass for the system lying midway between them:
This situation occurs commonly with binary stars (two stars bound gravitationally to
each other so that they revolve around their common center of mass). In many binary

star systems the masses of the two stars are similar and Newton's correction to Kepler's
3rd Law is very large.
Here is a Java applet that implements Newton's modified form of Kepler's 3rd law for
two objects (planets or stars) revolving around their common center of mass. By making
one mass much larger than the other in this interactive animation you can illustrate the
ideas discussed above and recover Kepler's original form of his 3rd Law where a less
massive object appears to revolve around a massive object fixed at one focus of an
ellipse.
These limiting cases for the location of the center of mass are perhaps familiar from our
afore-mentioned playground experience. If persons of equal weight are on a see-saw, the
fulcrum must be placed in the middle to balance, but if one person weighs much more
than the other person, the fulcrum must be placed close to the heavier person to achieve
balance.
Here is a Kepler's Laws Calculator that allows you to make simple calculations for
periods, separations, and masses for Keplers' laws as modified by Newton (see
subsequent section) to include the effect of the center of mass. (Caution: this applet is
written under Java 1.1, which is only supported by the most recent browsers. It should
work on Windows systems under Netscape 4.06 or the most recent version of Internet
Explorer 4.0, but may not yet work on Mac or Unix systems or earlier Windows
browsers.)
Weight and the Gravitational Force

We have seen that in the Universal Law of Gravitation the crucial quantity is mass. In
popular language mass and weight are often used to mean the same thing; in reality they
are related but quite different things. What we commonly call weight is really just the
gravitational force exerted on an object of a certain mass. We can illustrate by choosing
the Earth as one of the two masses in the previous illustration of the Law of Gravitation:

Thus, the weight of an object of mass m at the surface of the Earth is obtained by
multiplying the mass m by the acceleration due to gravity, g, at the surface of the Earth.
The acceleration due to gravity is approximately the product of the universal
gravitational constant G and the mass of the Earth M, divided by the radius of the Earth,
r, squared. (We assume the Earth to be spherical and neglect the radius of the object
relative to the radius of the Earth in this discussion.) The measured gravitational
acceleration at the Earth's surface is found to be about 980 cm/second/second.
Mass and Weight
Mass is a measure of how much material is in an object, but weight is a measure of the
gravitational force exerted on that material in a gravitational field; thus, mass and weight
are proportional to each other, with the acceleration due to gravity as the proportionality
constant. It follows that mass is constant for an object (actually this is not quite true, but
we will save that surprise for our later discussion of the Relativity Theory), but weight
depends on the location of the object. For example, if we transported the preceding
object of mass m to the surface of the Moon, the gravitational acceleration would change
because the radius and mass of the Moon both differ from those of the Earth. Thus, our
object has mass m both on the surface of the Earth and on the surface of the Moon, but it
will weigh much less on the surface of the Moon because the gravitational acceleration
there is a factor of 6 less than at the surface of the Earth.
Moment of Inertia
The moment of inertia of a solid body with density with respect to a given axis is
defined by the volume integral
(1)where is the perpendicular distance from the axis of rotation. This can be broken
into components as
(2)for a discrete distribution of mass, where r is the distance to a point (not the
perpendicular distance) and is the Kronecker delta, or
(3)for a continuous mass distribution. Depending on the context, I may be viewed either
as a tensor or a matrix. Expanding (3) in terms of Cartesian axes gives the equation
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(4)Similarly, the moment of inertia for a lamina with surface density is given by
(5)which can be written as
(6)The moment of inertia tensor I is symmetric, and is related to the angular momentum
vector L by
(7)where is the angular velocity vector.
The principal moments are given by the entries in the diagonalized moment of inertia
matrix, and are denoted (for a solid) A, B, and C in order of decreasing magnitude. In the
principal axes frame, the moments are also sometimes denoted , , and . The
principal axes of a rotating body are defined by finding values of I such that
(8)
(9)which is an eigenvalue problem.
The following table summarizes the moments of inertia of some common solids around
some of their principal axes.
solid/axis
moment of inertia
cylinder about symmetry axis
cylinder about central diameter
ellipsoid about principal axis
elliptical slab about major axis

elliptical slab about vertical
rectangular parallelepiped about major axis
ring about perpendicular axis
ring about diameter
rod about end
rod about center
sphere about diameter
spherical shell
torus about diameter
torus about symmetry axis
The following table summarizes the moments of inertia of some common laminae
around some of their principal axes.
lamina/axis
moment of inertia
ellipse about minor axis
rectangle about minor axis

Moment of Inertia and Angular Momentum
For linear, or translational, motion an object's resistance to a change in its state of motion
is called its inertia and is measured in terms of its mass, in kg. When a rigid body is
rotated, its resistance to a change in its state of rotation is called its rotational inertia, or
moment of inertia. This resistance has a two-fold property. (1) The amount of mass
present in the object and (2) the distribution of that mass about the chosen axis of
rotation. In general, the formula for an object's moment of inertia is CM = kmr2 where k
is a constant whose value varies from 0 to 1. Different positions of the axis result in
different moments of inertia for the same object; the further the mass is distributed from
the axis of rotation, the greater the value of its moment of inertia. Below is a series of
diagrams for a thin rod illustrating how the moment of inertia for the same object can
change with the placement of the axis of rotation.

axis on the far end of a thin rod:
I = 16/48 mL2
axis one-fourth of the way from the
end of a thin rod:
I = 7/48 mL2
axis at the center of a thin rod:
I = 4/48 mL2
These results would indicate that a thin rod would be most easily rotated about an axis
through its center of mass ( = 4/48 mr2) than about one of its far ends ( = 16/48 mr2).
Consider a majorette. If she twirls her baton about its center of gravity, for the same
amount of torque she will achieve a greater rate of angular acceleration than if she twirls
the same baton about a pivot closer to one of its ends. Whereas, a drum major
maneuvering his mace or a member of the band's flag corps spinning their flag poles
have work against a larger moment of inertia and therefore have more difficulty
accelerating their respective apparatus.
Three other shapes whose moments of inertia are important are: solid spheres, solid
disks and cylinders, and thin rings.

solid spheres
I = 2/5 mr2
solid disks and cylinders

I = 1/2 mr2
thin rings
I = mr2
Remember that the smaller the coefficient of mr 2, the easier it is to accelerate the object.
That is, spheres accelerate easier than cylinders, which accelerate easier than thin rings
or hoops. Since an object's moment of inertia increases as its mass is moved further from
its axis of rotation, hoops and rings would represent the greatest inertia in this collection
since all of their mass is concentrated at a constant distance, r, from the center of
rotation.
Based on our discussion, which object in the following picture would you predict would
reach the bottom of the incline first: the solid cylinder or the thin ring? To decide,
consider the rotational inertia of each object and how inertia affects motion.
Both objects have the same mass and equal diameters.
Angular momentum is the product of an object's moment of inertia (it rotational mass)
and its angular velocity. Angular momentum is a vector quantity represented by the
variable, L.
L= 
The units for angular momentum are: (kg m2)(radians/sec) = kg m2/sec. Note that
although the angular velocity must be expressed in radians/sec, the term radian is

dropped when expressing the units for its angular momentum. Remember that  =
2f. This expression was first introduced when we studied the sinusoidal equations for
SHM.
The vector nature of L is determined by the right hand rule. When your fingers curl in
the direction of the object's angular velocity, your thumb points in the direction of the
object's angular momentum. Examining the rotating masses illustrated in the diagrams
above, the sphere, disk and cylinder have angular velocities producing angular
momentum vectors pointing along the negative y-axis. The angular velocity of the thin
ring results in its angular momentum pointing along the positive x-axis.
Often we are required to determine the angular momentum of a point mass. A few
examples of point masses would be: (1) a speck of dust on a spinning CD's surface; (2) a
stopper moving in a circle at the end of a string; (3) a planet or asteroid moving in
circular orbit about the sun. As you can see, a point mass comes in many sizes. The term
applies to the fact that all of an object's mass is constrained to a small radius in
comparison to the radius of its circular motion or from the pivot point of the system.
That is, it can easily be represented by a single concentration of mass at the object's
center of mass.
The moment of inertia for a point mass traveling in a circle is = mr² and the
instantaneous tangential velocity of a point mass, v, equals rThis relationship
between angular and linear velocities can be understood by imagining a rotating
platform.
All points on the platform share the same angular velocity (they all pass through the
same angular displacement in a stated interval of time), but each one has a unique linear,
or tangential, velocity based on how far it is located from the axis of rotation .... that is,
how large a circumference it must travel through during each revolution. For the three
horse figurines shown above, the figurine closest to the central axis would have the least
tangential velocity since the fraction of its circumference that it travels during the 1/6 th
cycle shown is the shortest.

For a point mass, its angular momentum can be expressed as
L =  where = mr2 and  = v/r
L = (mr2)(v/r)
Lpoint mass = mvr
Note that this expression is the cross product of the object's linear momentum, mv, with
its radius, r. That is, the angular momentum of a point mass equals the product of the
magnitude of its linear momentum times its moment arm - the perpendicular distance
from the line of action of the momentum (instantaneous velocity vector) to the central
pivot or axis of rotation.
Law of Conservation of Angular Momentum

1 1 = 2 2
Angular momentum is conserved whenever there is no external force exerting a torque

on the object. An example of this occurs in skating. A skater spinning with arms out has
a greater , but a smaller angular velocity compared to when she is spinning with her
arms folded in ( is small, angular velocity is large).
Conservation of angular momentum justifies the relationship shown in Kepler's Second

Law,
A line from the planet to the sun sweeps out equal areas of space in equal intervals of
time.

The radius of gyration (r)
The radius of gyration (r) describes the way in which the area of a cross-section is
distributed around its centroidal axis. If the area is concentrated far from the centroidal
axis it will have a greater value of r and a greater resistance to buckling. A cross-section
can have more than one radius of gyration and most sections have at least two. If this is
the case, the section tends to buckle around the axis with the smallest value. The radius
of gyration is defined as:
r = sqrt (I/A)
where
r = radius of gyration
I = moment of inertia
A = area of the cross section
All things being equal, a circular pipe is the most efficient column section to resist
buckling. This is because it has an equal radius of gyration in all directions and it has the
its area distributed as far away as possible from the centroid.
Radius of gyration
The radius of gyration of an area with respect to a particular axis is the square root of
the quotient of the moment of inertia divided by the area. It is the distance at which the
entire area must be assumed to be concentrated in order that the product of the area and
the square of this distance will equal the moment of inertia of the actual area about the
given axis.
In other words, the radius of gyration describes the way in which the total cross-sectional
area is distributed around its centroidal axis. If more area is distributed further from the
axis, it will have greater resistance to buckling. The most efficient column section to
resist buckling is a circular pipe, because it has its area distributed as far away as
possible from the centroid.
If a cross-section has more than one radius of gyration (it is not circular) it will tend to
buckle around the axis with the smallest value.
The numerical value of the radius of gyration, r, is given by the following formula in
which I is the moment of inertia and A is the area:

The radius of gyration of a mass is similar except that

the moment of inertia of the mass is involved. The
numerical value of the radius of gyration of a mass, k,
is given by the following formula in which I is the
mass moment of inertia and m is the mass
Work power and energy
Definition and Mathematics of Work
In the first three units of The Physics Classroom, we utilized Newton's laws to analyze
the motion of objects. Force and mass information were used to determine the
acceleration of an object. Acceleration information was subsequently used to determine
information about the velocity or displacement of an object after a given period of time.
In this manner, Newton's laws serve as a useful scheme for analyzing motion and
making predictions about the final state of an object's motion. In this unit, an entirely
different scheme will be utilized to analyze the motion of objects. Motion will be
approached from the perspective of work and energy. The effect that work has upon the
energy of an object (or system of objects) will be investigated; the resulting velocity
and/or height of the object can then be predicted from energy information. In order to
understand this work-energy approach to the analysis of motion, it is important to first
have a solid understanding of a few basic terms. Thus, Lesson 1 of this unit will focus on
the definitions and meanings of such terms as work, mechanical energy, potential
energy, kinetic energy, and power.
In physics, work is defined as a force acting upon an object to cause a displacement.

There are three key words in this definition - force, displacement, and cause. In order for
a force to qualify as having done work on an object, there must be a displacement and
the force must cause the displacement. There are several good examples of work which
can be observed in everyday life - a horse pulling a plow through the fields, a father
pushing a grocery cart down the aisle of a grocery store, a freshman lifting a backpack
full of books upon her shoulder, a weightlifter lifting a barbell above her head, a shot-put
launching the shot, etc. In each case described here there is a force exerted upon an
object to cause that object to be displaced.

Read the following five statements and determine whether or not they represent
examples of work. Then depress the mouse upon the pop-up menu to view the answers.
Statement
Answer with Explanation
A teacher applies a force to a wall and becomes exhausted.
A book falls off a table and free falls to the ground.
A waiter carries a tray full of meals above his head by one arm across the room.
(Careful! This is a very difficult question which will be discussed in more detail later.)
A rocket accelerates through space.
Mathematically, work can be expressed by the following equation.
where F = force, d = displacement, and the angle (theta) is defined as the angle between
the force and the displacement vector. Perhaps the most difficult aspect of the above
equation is the angle "theta." The angle is not just any 'ole angle, but rather a very
specific angle. The angle measure is defined as the angle between the force and the
displacement. To gather an idea of its meaning, consider the following three scenarios.
 Scenario A: A force acts rightward upon an object as it is displaced rightward. In

such an instance, the force vector and the displacement vector are in the same
direction. Thus, the angle between F and d is 0 degrees.
 Scenario B: A force acts leftward upon an object which is displaced rightward. In
such an instance, the force vector and the displacement vector are in the opposite
direction. Thus, the angle between F and d is 180 degrees.
 Scenario C: A force acts upward upon an object as it is displaced rightward. In
such an instance, the force vector and the displacement vector are at right angles
to each other. Thus, the angle between F and d is 90 degrees.
Let's consider Scenario C above in more detail. Scenario C involves a situation similar to
the waiter who carried a tray full of meals above his head by one arm across the room. It
was mentioned earlier that the waiter does not do work upon the tray as he carries it
across the room. The force supplied by the waiter on the tray is an upward force and the
displacement of the tray is a horizontal displacement. As such, the angle between the
force and the displacement is 90 degrees. If the work done by the waiter on the tray were
to be calculated, then the results would be 0.
Regardless of the magnitude of the force and displacement, F*d*cosine 90 degrees is 0

(since the cosine of 90 degrees is 0). A vertical force can never cause a horizontal
displacement; thus, a vertical force does not do work on a horizontally displaced object!!

The equation for work lists three variables - each variable is associated with the one of
the three key words mentioned in the definition of work (force, displacement, and
cause). The angle theta in the equation is associated with the amount of force which
causes a displacement. As mentioned in a previous unit, when a force is exerted on an
object at an angle to the horizontal, only a part of the force contributes to (or causes) a
horizontal displacement. Let's consider the force of a chain pulling upwards and
rightwards upon Fido in order to drag Fido to the right. It is only the horizontal
component of the tensional force in the chain which causes Fido to be displaced to the
right. The horizontal component is found by multiplying the force F by the cosine of the
angle between F and d. In this sense, the cosine theta in the work equation relates to the
cause factor - it selects the portion of the force which actually causes a displacement.
When determining the measure of the angle in the work equation, it is important to
recognize that the angle has a precise definition - it is the angle between the force and
the displacement vector. Be sure to avoid mindlessly using any 'ole angle in the equation.
For instance, consider the activity performed in the "It's All Uphill" lab. A force was
applied to a cart to pull it up an incline at constant speed. Several incline angles were
used; yet, the force was always applied parallel to the incline. The displacement of the
cart was also parallel to the incline. Since F and d were in the same direction, the angle
was 0 degrees. Nonetheless, most students experienced the strong temptation to measure
the angle of incline and use it in the equation. Don't forget: the angle in the equation is
not just any 'ole angle; it is defined as the angle between the force and the displacement
vector.
Whenever a new quantity is introduced in physics, the standard metric units associated
with that quantity are discussed. In the case of work (and also energy), the standard
metric unit is the Joule (abbreviated "J"). One Joule is equivalent to one Newton of force
causing a displacement of one meter. In other words,
The Joule is the unit of work.

1 Joule = 1 Newton * 1 meter
1J = 1 N * m
In fact, any unit of force times any unit of displacement is equivalent to a unit of work.
Some nonstandard units for work are shown below. Notice that when analyzed, each set
of units is equivalent to a force unit times a displacement unit.
In summary, work is a force acting upon an object to cause a displacement. When a force
acts to cause an object to be displaced, three quantities must be known in order to
calculate the amount of work. Those three quantities are force, displacement and the
angle between the force and the displacement

Calculating the Amount of Work Done by Forces
In a previous part of Lesson 1, work was defined as a force acting upon an object to
cause a displacement. When a force acts to cause an object to be displaced, three
quantities must be known in order to calculate the work. Those three quantities are force,
displacement and the angle between the force and the displacement. The work is
subsequently calculated as force*displacement*cosine(theta) where theta is the angle
between the force and the displacement vectors. In this part of Lesson 1, the concepts
and mathematics of work will be applied in order to analyze a variety of physical
situations.
Express your understanding of the concept and mathematics of work by answering the
following questions. Depress the mouse on the "pop-up" menu to view the answers.
1. Apply the work equation to determine the amount of work done by the applied force
in each of the three situations described below.
2. On many occasions, there is more than one force acting upon an object. A free-body
diagram is a diagram which depicts the type and the direction of all the forces acting
upon an object. The following descriptions and their accompanying free-body diagrams
show the forces acting upon an object. For each case, indicate which force(s) are doing
work upon the object. Then calculate the work done by these forces. Use the "pop-up"
menu to view the answers.
Free-Body Diagram
A 10-N forces is applied to push a block across a friction free surface for a
displacement of 5.0 m to the right.

A 10-N frictional force slows a moving block to a stop after a displacement of 5.0 m to

the right.
A 10-N force is applied to push a block across a frictional surface at constant speed for
a displacement of 5.0 m to the right.

A 2-kg object is sliding at constant speed across a friction free surface for a
displacement of 5 m to the right.

A 2-kg object is pulled upward at constant speed by a 20-N force for a vertical
displacement of 5 m.
3. Before beginning its initial descent, a roller coaster car is always pulled up the first
hill to a high initial height. Work is done on the car (usually by a chain) to achieve this
initial height. A coaster designer is considering three different incline angles at which to
drag the 2000-kg car train to the top of the 60-meter high hill. In each case, the force
applied to the car will be applied parallel to the hill. Her critical question is: which angle
would require the most work? Analyze the data, determine the work done in each case,
and answer this critical question.
4. Ben Travlun carries a 200-N suitcase up three flights of stairs (a height of 10.0 m) and
then pushes it with a horizontal force of 50.0 N at a constant speed of 0.5 m/s for a
horizontal distance of 35.0 meters. How much work does Ben do on his suitcase during
this entire motion?
5. A force of 50 N acts on the block at the angle shown in the diagram. The block moves
a horizontal distance of 3.0 m.

6. How much work is done by an applied force to lift a 15-Newton block 3.0 meters
vertically at a constant speed?
7. A student with a mass of 80.0 kg runs up three flights of stairs in 12.0 sec. The student
has gone a vertical distance of 8.0 m. Determine the amount of work done by the student
to elevate his body to this height. Assume that her speed is constant.
Potential Energy
An object can store energy as the result of its position. For example, the heavy ram of a
pile driver is storing energy when it is held at an elevated position. This stored energy of
position is referred to as potential energy. Similarly, a drawn bow is able to store energy
as the result of its position. When assuming its usual position (i.e., when not drawn),
there is no energy stored in the bow. Yet when its position is altered from its usual
equilibrium position, the bow is able to store energy by virtue of its position. This stored
energy of position is referred to as potential energy. Potential energy is the stored energy
of position possessed by an object.
The two examples above illustrate the two forms of potential energy to be discussed in
this course - gravitational potential energy and elastic potential energy. Gravitational
potential energy is the energy stored in an object as the result of its vertical position (i.e.,
height). The energy is stored as the result of the gravitational attraction of the Earth for
the object. The gravitational potential energy of the heavy ram of a pile driver is
dependent on two variables - the mass of the ram and the height to which it is raised.
There is a direct relation between gravitational potential energy and the mass of an
object; more massive objects have greater gravitational potential energy. There is also a

direct relation between gravitational potential energy and the height of an object; the
higher that an object is elevated, the greater the gravitational potential energy. These
relationships are expressed by the following equation:
PEgrav = mass * g * height

PEgrav = m * g * h
In the above equation, m represents the mass of the object, h represents the height of the
object and g represents the acceleration of gravity (approximately 10 m/s/s on Earth).
To determine the gravitational potential energy of an object, a zero height position must
first be arbitrarily assigned. Typically, the ground is considered to be a position of zero
height. But this is merely an arbitrarily assigned position which most people agree upon.
Since many of our labs are done on tabletops, it is often customary to assign the tabletop
to be the zero height position; again this is merely arbitrary. If the tabletop is the zero
position, then the potential energy of an object is based upon its height relative to the
tabletop. For example, a pendulum bob swinging to and from above the table top has a
potential energy which can be measured based on its height above the tabletop. By
measuring the mass of the bob and the height of the bob above the tabletop, the potential
energy of the bob can be determined.
Since the gravitational potential energy of an object is directly proportional to its height
above the zero position, a doubling of the height will result in a doubling of the
gravitational potential energy. A tripling of the height will result in a tripling of the
gravitational potential energy. Use this principle to determine the blanks in the following
diagram. Knowing that the potential energy at the top of the tall pillar is 30 J, what is the
potential energy at the other positions shown on the hill and the stairs.

The second form of potential energy which we will discuss in this course is elastic
potential energy. Elastic potential energy is the energy stored in elastic materials as the
result of their stretching or compressing. Elastic potential energy can be stored in rubber
bands, bungee chords, trampolines, springs, an arrow drawn into a bow, etc. The amount
of elastic potential energy stored in such a device is related to the amount of stretch of
the device - the more stretch, the more stored energy.
Springs are a special instance of a device which can store elastic potential energy due to
either compression or stretching. A force is required to compress a spring; the more
compression there is, the more force which is required to compress it further. For certain
springs, the amount of force is directly proportional to the amount of stretch or
compression (x); the constant of proportionality is known as the spring constant (k).
Such springs are said to follow Hooke's Law. If a spring is not stretched or compressed,
then there is no elastic potential energy stored in it. The spring is said to be at its
equilibrium position. The equilibrium position is the position that the spring naturally
assumes when there is no force applied to it. In terms of potential energy, the equilibrium
position could be called the zero-potential energy position. There is a special equation
for springs which relates the amount of elastic potential energy to the amount of stretch
(or compression) and the spring constant. The equation is

To summarize, potential energy is the energy which an object has stored due to its
position relative to some zero position. An object possesses gravitational potential
energy if it is positioned at a height above (or below) the zero height position. An object
possesses elastic potential energy if it is at a position on an elastic medium other than the
equilibrium position.
Check Your Understanding
Check your understanding of the concept of potential energy by answering the following
questions. When finished, use the pop-up menu to view the answers.
1. A cart is loaded with a brick and pulled at constant speed along an inclined plane to
the height of a seat-top. If the mass of the loaded cart is 3.0 kg and the height of the seat
top is 0.45 meters, then what is the potential energy of the loaded cart at the height of the
seat-top?
2. If a force of 15.0 N is used to drag the loaded cart (from previous question) along the
incline for a distance of 0.90 meters, then how much work is done on the loaded cart?
Note that the work done to lift the loaded cart up the inclined plane at constant speed is
equal to the potential energy change of the cart. This is not coincidental! The reason for
the relation between the potential energy change of the cart and the work done upon it is
the subject of Lesson
Kinetic Energy
Kinetic energy is the energy of motion. An object which has motion - whether it be
vertical or horizontal motion - has kinetic energy. There are many forms of kinetic
energy - vibrational (the energy due to vibrational motion), rotational (the energy due to
rotational motion), and translational (the energy due to motion from one location to
another). To keep matters simple, we will focus upon translational kinetic energy. The
amount of translational kinetic energy (from here on, the phrase kinetic energy will refer
to translational kinetic energy) which an object has depends upon two variables: the
mass (m) of the object and the speed (v) of the object. The following equation is used to
represent the kinetic energy (KE) of an object.

where m = mass of object
v = speed of object
his equation reveals that the kinetic energy of an object is directly proportional to the
square of its speed. That means that for a twofold increase in speed, the kinetic energy
will increase by a factor of four; for a threefold increase in speed, the kinetic energy will
increase by a factor of nine; and for a fourfold increase in speed, the kinetic energy will
increase by a factor of sixteen. The kinetic energy is dependent upon the square of the
speed. As it is often said, an equation is not merely a recipe for algebraic problem-
solving, but also a guide to thinking about the relationship between quantities.
Kinetic energy is a scalar quantity; it does not have a direction. Unlike velocity,
acceleration, force, and momentum, the kinetic energy of an object is completely
described by magnitude alone. Like work and potential energy, the standard metric units
of measurement for kinetic energy is the Joule. As might be implied by the above
equation, 1 Joule is equivalent to 1 kg*(m/s)^2.
Use your understanding of kinetic energy to answer the following questions. Then
depress the mouse on the "pop-up menu" to view the answers.
1. Determine the kinetic energy of a 1000-kg roller coaster car that is moving with a
speed of 20.0 m/s.
2. If the roller coaster car in the above problem were moving with twice the speed, then
what would be its new kinetic energy?
3. Missy Diwater, the former platform diver for the Ringling Brother's Circus had a
kinetic energy of 15 000 J just prior to hitting the bucket of water. If Missy's mass is 50
kg, then what is her speed?
4. A 750-kg compact car moving at 100 km/hr has approximately 290 000 Joules of
kinetic energy. What is the kinetic energy of the same car if it is moving at 50 km/hr?
(HINT: use the kinetic energy equation as a "guide to thinking.")

Mechanical Energy
In a previous part of Lesson 1, it was said that work is done upon an object whenever a
force acts upon it to cause it to be displaced. Work is a force acting upon an object to
cause a displacement. In all instances in which work is done, there is an object which
supplies the force in order to do the work. If a World Civilization book is lifted to the top
shelf of a student locker, then the student supplies the force to do the work on the book.
If a plow is displaced across a field, then some form of farm equipment (usually a tractor
or a horse) supplies the force to do the work on the plow. If a pitcher winds up and
accelerates a baseball towards home plate, then the pitcher supplies the force to do the
work on the baseball. If a roller coaster car is displaced from ground level to the top of
the first drop of the Shock Wave, then a chain (driven by a motor) supplies the force to
do the work on the car. If a barbell is displaced from ground level to a height above a
weightlifter's head, then the weightlifter is supplying a force to do work on the barbell.
In all instances, an object which possesses some form of energy supplies the force to do
the work. In the instances described here, the objects doing the work (a student, a tractor,
a pitcher, a motor/chain) possess chemical potential energy stored in food or fuel which
is transformed into work. In the process of doing work, the objects doing the work
exchange energy in one form to do work on another object to give it energy. The energy
acquired by the objects upon which work is done is known as mechanical energy.
Mechanical energy is the energy which is possessed by an object due to its motion or its
stored energy of position. Mechanical energy can be either kinetic energy (energy of
motion) or potential energy (stored energy of position). Objects have mechanical energy
if they are in motion and/or if they are at some position relative to a zero potential
energy position (for example, a brick held at a vertical position above the ground or zero
height position). A moving car possesses mechanical energy due to its motion (kinetic
energy). A moving baseball possesses mechanical energy due to both its high speed
(kinetic energy) and its vertical position above the ground (gravitational potential
energy). A World Civilization book at rest on the top shelf of a locker possesses
mechanical energy due to its vertical position above the ground (gravitational potential
energy). A barbell lifted high above a weightlifter's head possesses mechanical energy
due to its vertical position above the ground (gravitational potential energy). A drawn
bow possesses mechanical energy due to its stretched position (elastic potential energy).
An object which possesses mechanical energy is able to do work. In fact, mechanical

energy is often defined as the ability to do work. Any object which possesses mechanical
energy - whether it be in the form of potential energy or kinetic energy - is able to do
work. That is, its mechanical energy enables that object to apply a force to another object
in order to cause it to be displaced.
Numerous examples can be given of how an object with mechanical energy can harness
that energy in order to apply a force to cause another object to be displaced. A classic
example involves the heavy ram of a pile driver. A pile driver consists of a massive
object which is elevated to a high position and allowed to fall upon another object

(called the pile) in order to drive it downwards. Upon hitting the pile, the ram applies a
force to it in order to cause it to be displaced. The diagram below depicts the process by
which the mechanical energy of a pile driver can be used to do work.
A hammer is a miniature version of a pile driver. The mechanical energy of a hammer

gives the hammer its ability to apply a force to a nail in order to cause it to be displaced.
Because the hammer has mechanical energy (in the form of kinetic energy), it is able to
do work on the nail. Mechanical energy is the ability to do work.
Another example which illustrates how mechanical energy is the ability of an object to
do work can be seen any evening at your local bowling alley. The mechanical energy of
a bowling ball gives the ball the ability to apply a force to a bowling pin in order to
cause it to be displaced. Because the massive ball has mechanical energy (in the form of
kinetic energy), it is able to do work on the pin. Mechanical energy is the ability to do
work.
A dart gun is still another example of how mechanical energy of an object can do work
on another object. When a dart gun is loaded and the springs are compressed, it
possesses mechanical energy. The mechanical energy of the compressed springs give the
springs the ability to apply a force to the dart in order to cause it to be displaced.
Because of the springs have mechanical energy (in the form of elastic potential energy),
it is able to do work on the dart. Mechanical energy is the ability to do work.
A common scene in the western United States is a "wind farm." High speed winds are
used to do work on the blades of a turbine at the so-called wind farm. The mechanical
energy of the moving air give the air the ability to apply a force and cause a
displacement of the blades. As the blades spin, their energy is subsequently converted
into electrical energy (a non-mechanical form of energy) and supplied to homes and
industries in order to run electrical appliances. Because the moving wind has mechanical
energy (in the form of kinetic energy), it is able to do work on the blades. Once more,
mechanical energy is the ability to do work.
As already mentioned, the mechanical energy of an object can be the result of its motion
(i.e., kinetic energy) and/or the result of its stored energy of position (i.e., potential
energy). The total amount of mechanical energy is merely the sum of the potential

energy and the kinetic energy. This sum is simply referred to as the total mechanical
energy (abbreviated TME).
TME = PE + KE
As discussed earlier, there are two forms of potential energy discussed in our course -
gravitational potential energy and elastic potential energy. Given this fact, the above
equation can be rewritten:
TME = PEgrav + PEspring + KE
The diagram below depicts the motion of Li Ping Phar (esteemed Chinese ski jumper) as
she glides down the hill and makes one of her record-setting jumps.
The total mechanical energy of Li Ping Phar is the sum of the potential and kinetic
energies. The two forms of energy sum up to 50 000 Joules. Notice also that the total
mechanical energy of Li Ping Phar is a constant value throughout her motion. There are

conditions under which the total mechanical energy will be a constant alue and
conditions under which it is a changing value. This is the subject of Lesson 2 - the work-
energy theorem. For now, merely remember that total mechanical energy is the energy
possessed by an object due to either its motion or its stored energy of position. The total
amount of mechanical energy is merely the sum of these two forms of energy. And
finally, an object with mechanical energy is able to do work on another object.
Power
The quantity work has to do with a force causing a displacement. Work has nothing to do
with the amount of time that this force acts to cause the displacement. Sometimes, the
work is done very quickly and other times the work is done rather slowly. For example, a
rock climber takes an abnormally long time to elevate her body up a few meters along
the side of a cliff. On the other hand, a trail hiker (who selects the easier path up the
mountain) might elevate her body a few meters in a short amount of time. The two
people might do the same amount of work, yet the hiker does the work in considerably
less time than the rock climber. The quantity which has to do with the rate at which a
certain amount of work is done is known as the power. The hiker has a greater power
rating than the rock climber.
Power is the rate at which work is done. It is the work/time ratio. Mathematically, it is
computed using the following equation.
The standard metric unit of power is the Watt. As is implied by the equation for power, a
unit of power is equivalent to a unit of work divided by a unit of time. Thus, a Watt is
equivalent to a Joule/second. For historical reasons, the horsepower is occasionally used
to describe the power delivered by a machine. One horsepower is equivalent to
approximately 750 Watts.
Most machines are designed and built to do work on objects. All machines are typically
described by a power rating. The power rating indicates the rate at which that machine
can do work upon other objects. Thus, the power of a machine is the work/time ratio for
that particular machine. A car engine is an example of a machine which is given a power
rating. The power rating relates to how rapidly the car can accelerate the car. Suppose
that a 40-horsepower engine could accelerate the car from 0 mi/hr to 60 mi/hr in 16
seconds. If this were the case, then a car with four times the horsepower could do the
same amount of work in one-fourth the time. That is, a 160-horsepower engine could
accelerate the same car from 0 mi/hr to 60 mi/hr in 4 seconds. The point is that for the
same amount of work, power and time are inversely proportional. Equations can be
"guides to thinking" about the relationships between quantities. The power equation
suggests that a more powerful engine can do the same amount of work in less time.

A person is also a machine which has a power rating. Some people are more power-full
than others; that is, they are capable of doing the same amount of work in less time or
more work in the same amount of time. In the Personal Power lab, students determined
their own personal power by doing work on their bodies to elevate it up a flight of stairs.
By measuring the force, displacement and time, we were able to measure our personal
power rating. Suppose that Ben Pumpiniron elevates his 80-kg body up the 2.0 meter
stairwell in 1.8 seconds. If this were the case, then we could calculate Ben's power
rating. It can be assumed that Ben must apply a 800-Newton downward force upon the
stairs to elevate his body. By so doing, the stairs would push upward on Ben's body with
just enough force to lift his body up the stairs. It can also be assumed that the angle
between the force of the stairs on Ben and Ben's displacement is 0 degrees. With these
two approximations, Ben's power rating could be determined as shown below.
Ben's power rating is 889 Watts; what a "horse."
The expression for power is work/time. Now since the expression for work is
force*displacement, the expression for power can be rewritten as
(force*displacement)/time. Yet since the expression for velocity is displacement/time,
the expression for power can be rewritten once more as force*velocity. This is shown
below.
This new expression for power reveals that a powerful machine is both strong (big force)
and fast (big velocity). The powerful car engine is strong and fast. The powerful farm
equipment is strong and fast. The powerful weightlifters are strong and fast. The
powerful linemen on a football team are strong and fast. A machine which is strong
enough to apply a big force to cause a displacement in a small mount of time (i.e., a big
velocity) is a powerful machine.
Use your understanding of work and power to answer the following questions. When
finished, depress the mouse on the "pop-up menu" to view the answers.

1. Two physics students, Will N. Andable and Ben Pumpiniron, are in the weightlifting
room. Will lifts the 100-pound barbell over his head 10 times in one minute; Ben lifts the
100-pound barbell over his head 10 times in 10 seconds. Which student does the most
work? Which student delivers the most power? Explain your answers.
2. During the Personal Power lab, Jack and Jill ran up the hill. Jack is twice as massive
as Jill; yet Jill ascended the same distance in half the time. Who did the most work? Who
delivered the most power? Explain your answers.
Work
Whenever a force displaces its point of application in a direction parallel to itself work is
done. Only a force can do work.
When work is done by a force, energy is transformed from one type to another.
e.g. 1: When a body falls accelerating under gravity, the work done is the amount by
which the potential energy of the system drops and its kinetic energy rises.
e.g. 2: When a body moves against a frictional force, the work done is the amount by
which sound and heat energy are produced at the expense of mechanical energy.
If the point of application of a constant force F is displaced by an amount s in a direction

making an angle with the force, then the work done is defined as:
Work is a scalar quantity.
When unit force displaces its point of application through unit distance in a direction
parallel to itself, unit work is done.
1 Joule = 1 newton x 1 metre [MKS]

1 erg = 1 dyne x 1 centimetre [cgs]
The joule and the erg are respectively the MKS and the cgs units of work. These units are
also called the absolute units of work as they are of constant value at all locations.
Another set of units of work - the gravitational units - are defined as the product of the

gravitational unit of force times the unit of distance.
A force does zero work if:
 there is no displacement of its point of application (s = 0)

e.g. No work is done by the force of static friction since there is no displacement
of its point of application
 the point of application is displaced perpendicular to itself (cos = 0)
e.g. There is no work done by the gravitational force (weight) when one walks on
a horizontal surface.
 the force itself is of negligible magnitude (F=0)
e.g. No work is done by intermolecular forces when a gas at low pressure
undergoes expansion.
Numericals
1.
Starting from rest an automobile of mass 500 kg accelerates at 1.5 m/s 2 for 12.0 second
and then continues to move at constant speed for 5 minutes. It is then stopped by a
constant retardation of 2 ms -2 . Assuming that the whole journey is performed on a level
road, the effect of friction being equivalent to a constant retarding force of 100 N,
calculate:
a. the total distance covered
b. the engine power used during motion at constant speed
c. during accelerated motion what is the
i. resultant force acting on the automobile
ii. average engine power used
[ 5589m;1800W;750N;7650W ]
2.
Refer to the graphs a] b] and c] given below, each of which shows variation of force F
with respect to displacement s in the direction of the force. For each case find the gain in
kinetic energy on covering a displacement of 4.0 metres.

3.
The effective length of a simple pendulum is 1.0 m and the mass of its bob is 40g. The
bob is displaced through an angle of 60o and released. Give a sketch showing the bob in
a displaced position and mark in all forces acting on the bob. Which of these forces does
work as the bob is released from this displaced position? Find the kinetic energy when
the string makes an angle of 30o with the vertical. [ 5.2 x 10-2 J ]
4.
A 12 kg block is pushed up through 20 m along sloping surface of a plane inclined at 37 o
to the horizontal by a constant force F of 120 N acting along the plane and parallel to the
displacement. If the frictional resistance is 24 N, find:
i. The work done by the applied force F
ii. The work done by the frictional force f
iii. The work done by the normal reaction
iv. The work done by the weight of the body
v. The increase in the potential energy of the block
vi. The increase in kinetic energy of the block
[ 2400J; -480J; 0J; -1440J; +1440 J; +480 J ]

5.
The system shown below is released from rest when the 12 kg block is 3.0 m above the
floor. Use the principle of conservation of energy to find the velocity with which the
block strikes the floor. The string is light and there is no friction between the pulley and
its axle.
[ 5.5 ms-1 ]
6.
Explain with reason whether the law of conservation of energy is violated in each of the
following cases:-
a. Raindrops fall through air with constant velocity.
b. When a bullet is fired from a gun, both the bullet and the gun gain in kinetic
energy starting from rest.
c. When one runs up and down stairs one gets exhausted even though the net
mechanical work done is zero in each repeat.
d. When a helium filled balloon is released, it rises upwards and moves faster at the
sametime.
Work-Energy relation
If there are opposing forces acting on a body then the work done ( W) by the net
resultant force F on a particle is equal to its gain in kinetic energy.
That is, W = K

where, K .= Rise in kinetic energy of the system

P .= Rise in potential energy of the system
U= Rise in thermal energy of the system
Energy of Freely Falling Body

Let a body of mass m initially stationary at a height h fall through a vertical distance x.
Let the acceleration due to gravity g remain constant through this small distance of fall.
Therefore, for a freely falling body as it loses potential energy it gains an equal amount
of kinetic energy if there is no other form of energy being produced (i.e. if there is no
frictional force). In case there is friction, then the body loses mechanical energy by the
same amount as is the energy dissipated as heat, sound etc.
Potential Energy
A body placed at a height above the earth and released subsequently under-goes
downward displacement i.e. towards its position of rest. The free end of a rubber cord
pulled taut and then released undergoes displacement towards its position of rest. A
simple pendulum displaced from its position of rest subsequently undergoes
displacement towards its rest position. All these bodies undergo displacement because of
their being placed away from their normal resting position. All these bodies therefore
posses energy by virtue of their displaced positions. Such energy is called potential
energy.
If a body has an external unbalanced force acting on it, only then can the body can move
itself if released. This external force is exerted upon the displaced body by another body
which in turn experiences an equal and opposite reaction force. Thus the two body
system has potential energy. A body may also have energy by virtue of the positions of
its multiple parts with respect to each other i.e.the energy is a property of a system of
many objects. This is its potential energy by virtue of its configuration.

The system comprising the apple and the earth has potential energy as there is relative
displacement (from the rest position) between its parts.
The system comprising cork, water and earth has potential energy by virtue of the
position of its various parts with respect to each other.
The system comprising the mass and the stretched spring has potential energy by virtue
of the position of its different parts with respect to each other.
The system comprising the two charged bodies has potential energy by virtue of the
position of its two parts with respect to each other.
If the displacement from rest position is oppositely directed to the system's force, then
the system's force does work by an amount

W = Fs s cos 180o = - FsS
This actually implies that an external force did positive work (displaced the body). The
system thus gains in potential energy by the amount of work done by the external force.
Work done by external force is thus given by:
= FE s cos 0o = + FES
[Since there is to be no gain in kinetic energy, there was no acceleration, and thus no
unbalance between the system's force and the external force acting. In magnitude, the
external force was equal to the system's force. The direction was opposite to the system's
force and thus the same as the displacement.]
If the displacement is in the direction of the system's force then the system's force does
positive work. the system loses this amount of potential energy. Increase in potential
energy is the work done by the external force and is given by the product of the external
force (equal in magnitude to the system's force) and the displacement parallel to the
force .
If the system's force FS is constant in magnitude and direction over the entire parallel
displacement s (=s' cos ), then the required external force FE is also a constant.
Thus, gain in potential energy = FE . s
In case the force varies with respect to the displacement, then the work done is given by
where FE (= FS ) is a function of the displacement s
[It is only in the case of a system comprising an object and the earth that the system's
force acting on the object is gravitational and equal to mg . It is only in this particular
case that a small displacement h in a vertical direction causes a change of the amount
mgh in its potential energy.
In the event that the vertical displacement for this system is large enough for the
acceleration due to gravity to alter, the potential energy differs from mg and is to be
found by integration as above
In case of the other systems discussed above the potential energy is not equal to mgh.]

CONSERVATION OF ENERGY
People have always tried to better understand the world around them. They have
described it, probed it, and theorized about it in an attempt to explain everyday
phenomena. There have been many theories about the world, many of which have been
wrong. For a very long time, man believed that the Earth was the center of the universe.
Aristotle believed that, when dropped from the same height, a heavier object would fall
to the ground faster than a lighter object. This idea was proven wrong more than 1000
years later when it was shown that, all objects, regardless of mass, fall to the ground with
the same acceleration.
People throughout the ages have come up with ways to describe the world with numbers
and theories. One idea that has come up over and over again is conservation: the idea
that certain quantities in nature always stay the same.
Numerous quantities have been found to be conserved. Three important conserved

quantities that we will be discussing are energy, linear momentum and angular
momentum.
There has been no known example of a situation in which energy, linear momentum or
angular momentum are not conserved. However, from a practical standpoint,
conservation ismost useful when we know how to calculate and count up exactlyhow
much of each quantity there is. We need to know how to count up how much energy,
linear momentum and angular momentum there is in a system in order to use
conservation to make predictions about the system.
Energy describes the nature of the system. Linear momentum is a quantity associated
with motion in a straight line and angular momentum is associated with motion about a
point. These three quantities are related to the state of the system. Because these
quantities are conserved, we can often use them to make both general and specific
predictions about how the system will behave in the future.
Conservation is all about accounting for the amounts of each quantity in every form and
at every location. The theory of conservation states that we should be able to
continuously account for every little bit of energy, linear momentum and angular
momentum; nothing just suddenly appears or disappears. In applying the theory, we are
limited by whether we can accurately count up what is there.
Often, we are only concerned with the changes in the quantities. By conservation,
everything must be accounted for; an increase in one form of energy means that there
was a decrease in another form. If one object gains linear momentum, another object
must have lost some. In cases such as these, we are comparing the quantities before and
after some interaction.

The whole idea behind conservation is that there can be no surprises. There is a source
for every change. We are sometimes unable to apply conservation because we can not
easily calculate the amount present of some quantity; for example, in some
circumstances we might have trouble calculating heat energy. However, this does not
mean that the energy is not conserved. Something else must have lost energy to produce
the heat energy and it just happens that we have a hard time counting up the amount of
heat energy directly. Conservation of energy, linear momentum and angular momentum
is always true; we just have to be careful in how we account for the quantities that we
know are conserved.
In application of momentum, the most important law to know is the Conservation of

Momentum. Momentum is said to be conserved when there is zero net external force
acting on the object.
To explain this mathematically:

 Definition of momentum: p = mv
 Newton's 1st law: F = ma => m dv/dt = dp/dt
 Net external force is zero: F = 0 => dv/dt = 0
 dv/dt = 0 => p = 0
Conservation of Momentum is mathematically expressed as:

mvi = mvf
When expanded to n-particle system the n-conservative momentum of each particle add:
Now let's take a look at some things to verify what we covered in previous sections.
In Aiming, a statement was made that the angle formed when the cue collides with the
object ball which is initially at rest is always 90 degrees. This is so if assuming there is
no friction and no rotation on the balls so that no kinetic energy is dissipated. This is the
reason why in real life (with friction and rotation) it is never perfectly right angular after
collision but very similar. Back to verifying why it is theoretically 90 degrees, we must
consider first the conservation of momentum and the conservation of kinetic energy:
All balls are equal in mass and the object ball is initially at rest, v2i = 0
Conservation of momentum gives us:

mv1i = mv1f + mv2f
v1i = v1f + v2f
This is a vector equation and geometrically means that three vectors form a triangle.
(1/2)mv1i2 = (1/2)mv1f2 + (1/2)mv2f2 v1i2 = v1f2 + v2f2

As you can see, the obtained velocities are in a form of Pythagorean theorem and the

vector velocity v1i is the hypotenuse of two leg vector velocities, v1f and v2f.
Hence, right angle between the two final vectors.
In Elasticity, we learned that the coefficient of restitution, e = -(v1i - v2i)/(v1f - v2f).

To verify this, we begin again by writing the conservation of kinetic energy:
(1/2)mv1i2 + (1/2)Mv2i2 = (1/2)mv1f2 + (1/2)Mv2f2
After rearrangement of terms:
mv1i2 - mv1f2 = Mv2f2 - Mv2i2
Eq1: m(v1i - v1f)(v1i + v1f) = M(v2f - v2i)(v2f + v2i)
To simplify the equation any further, we must look at two-object conservation of

momentum:
mv1i + Mv2i = mv1f + Mv2f
Eq2: m(v1i - v1f) = M(v2f - v2i)
By inspection, we notice that a part of Eq1 may be simplified using Eq2 and obtain:
v1i + v1f = v2i + v2f OR
v1i - v2i = -(v1f - v2f)
Eq3: -(v1i - v2i)/(v1f - v2f) = 1
We notice from Eq3 which is the coefficient of restitution for a perfectly elastic collision,
e = 1. This is true because in elastic collisions, the kinetic energy is conserved and hence
the linear momentum is also conserved.
Linear Momentum and Collisions
Seems like everything wants to keep on keeping on...

Now lets extend our understanding of translational motion by defining a new term,
"linear momentum". Linear momentum of a particle is the product of a particle's mass
and its velocity. For some reason which escapes me at the moment, the symbol p is
customarily used to designate momentum. Since it is the product of a vector and a scalar,
it is a vector quantity, p = m * v. The rate of change of momentum with respect to time,
taking mass to be constant, is p' = m * v', where v' is the vector acceleration of the
particle. Remember that we are using the ' symbol to represent the rate of change with
respect to time. From our expression of Newton's second law, we know that m * v' is the
force on a particle so p' = f. In fact when Newton published his second law it was in the
form p' = f, rather than as we first expressed it.
For a system of particles, the linear momentum is just the vector sum of that of each of
the particles individually. As we worked out in the previous lesson, m * vcm = m1*v1 +
m2*v2...+ mn*vn so for a system of particles p = m * vcm, and p' = m * acm = fa, where fa
is the force resultant external force applied to the particles.
In words this says that for a system of particles, the rate of change of the total
momentum is equal to the applied force. Run the Two Particle Center of Mass display
again and look at the center of mass track. Since the objects are moving with constant
velocity we know that there is no force applied to the two object system. The momentum

of the whole ystem must be constant, meaning that the center of mass travels in a straight
line at constant velocity.
This idea that in the absence of external forces the momentum remains constant is called
the principle of the "conservation of linear momentum". This principle, like the
conservation of energy or work-energy theorem, seems to be near the heart of the laws of
nature. It applies in all kinds of situations, even those outside the realm of classical
dynamics.
So far we have developed two distinct modeling schemes for building mathematical
models of dynamical systems. In the first instance we used Newton's second law. The
implicit assumption in using that modeling scheme is that we know what forces are
applied to the particle and how that force varies as a function of time. The other
modeling scheme involved the first fundamental conservation law we discovered, the
conservation of mechanical energy. In applying that model we need to know the way the
potential energy varies as a function of position. Now that we have the conservation of
linear momentum as another basic principle we can model some situations in which we
know neither the forces nor the potential field.
One instance of a system where the detail of forces and fields is obscured is in collisions.
It is theoretically possible to know in detail the way force varies with time in a collision
but as a practical matter that may be very difficult to determine. By "collision" , I mean
an interaction between two objects which is over is a very short time, compared to the
time of observation. With this as a working definition, we may divide time up into three
regions, before collision, during collision and after collision. Using the conservation of
linear momentum and the conservation of mechanical energy, we may determine from
the before collision conditions what the conditions will be after the collision without
ever knowing exactly what went on during the collision.
You may recall one example of a collision already illustrated in the Billiards display. It
was used in the discussion of Newton's third law. I cheated a little bit and used the
conservation of linear momentum in that display before we understood the concept.
Perhaps you would like to review that display with the idea of conservation of
momentum in mind.
The model behind this display operates on Newton's first law during the times before
and after the collision. We assume there are no forces on the balls before and after the
collision so the equations of motion are the trivial case of no change in velocity. During
the collision the two balls compress one another slightly exerting huge forces. Even
though we will not work out in detail the variation of force with time, we can get some
notion of the total force involved.
We know now that f=p'. Lets estimate the forces involved in a head on, dead center
collision between the two balls of equal mass. The cue ball stops in that case and the red
ball goes off at the velocity that the cue ball had initially. Verify this with the model or
go down to your local tavern and try it. The way to get a head on collision with the

model is to place the cursor right on the center of the red ball for the shot. You might
want to start with the cue ball at the far right for maximum shot speed. Let's assume you
achieve a cue ball speed of 2 meters per second. The change in velocity of the cue ball
then is 2m/s since it comes to a complete stop.
A normal cue ball has a mass of about .25kg so the momentum change is .25*2kgm/s =
0.5kgm/s. The question now becomes how long it takes for the cue ball to stop. Here is
where we have to estimate a number. Let's say the balls remain in contact for 1e-3
seconds. The average force magnitude |f| then during the 1e-3 seconds is |f| = p' = .5/1e-3
= 500 Newtons
This is considerable force, perhaps enough to lift a small adult off the ground. Forces of
the sort we are talking about here, large but short are called "impulse" forces. Impulse
forces even though of short duration can do considerable damage to the objects involved.
In general they are estimated by dividing a change in momentum by the time to make the
change. Conversely the average force if known, multiplied by the time the collision lasts
is called the impulse of the force and is equal to the change in momentum.
The basis for our collision model will be the two conservation principles, conservation
of kinetic energy and conservation of linear momentum. We can work with conservation
of kinetic energy only, because in studying collisions we want to isolate the effect of the
collision from the effect of other forces. So we will set up our conditions such that other
forces like those arising from potential energy changes do not apply. Actually collisions
are classified as to whether or not kinetic energy is conserved. If it is, the collisions are
called "elastic" collisions, if not they are called inelastic collisions. We are going to limit
ourselves to elastic collisions.
Consider two spherical objects one of mass m1 and the other of mass m2. Let's set things
up so these objects are approaching each other along the line joining their centers, a
recipe for a head on collision. Let these objects not be subject to any forces and be not
rotating or vibrating. The motion is purely translation. Under these conditions we may
choose a reference frame which is one dimensional, with the objects on the x axis.
Applying the conservation of linear momentum to this situation we have, m1 * v1i + m2 *

v2i = m1 * v1 + m2 * v2, where vi is the initial velocity of each object and v is the final
velocity of each object. The conservation of kinetic energy gives us the equation 1/2 *
m1 * v1i2 + 1/2 * m2 * v2i2 = 1/2 * m1 * v12 + 1/2 * m2 * v22. These may be rewritten as
follows:
m1 * (v1i - v1) = m2 * (v2 - v2i) for the momentum equation and m1 * (v1i2 - v12) = m2 * (v22
- v2i2) for the energy.
As long as the difference between final and initial velocities is not zero for either object
(meaning a collision actually happens), we may divide the second equation by the first
one which yields v1i + v1 = v2 + v2i , or v1i - v2i = v2 - v1. In other words in a one
dimensional elastic collision, the relative velocity of approach before the collision equals
the relative velocity of separation after collision.

Now here is one of those places where we are going to lean on that algebra background I
told you would need in this course. To get the final velocities in terms of the initial
velocities and the masses, you would solve the last equation above for v2 and plug that
into the momentum equation and solve to get v1 = v1i * (m1 - m2) / (m1 + m2) + v2i * (2 *
m2) / (m1 + m2). Likewise v2 = v1i * (2 * m1) / (m1 + m2) + v2i * (m2 - m1) / (m2 + m1).
So we have the basis for a model of a one-dimensional elastic collision. For initial
conditions v1i and v2i, if a collision happens, the final velocities depend on the masses as
above. In the special case where the objects have the same mass and object 2 was
initially at rest, you can easily see from the equations that object 1 comes to a stop and
object 2 takes off with the same velocity object 1 originally had. The objects just
exchange velocities.
One-dimensional (head on) elastic collisions though leave a lot of naturally occurring
collisions unaccounted for. We should extend this model to two dimensions to be more
realistic. In three dimensions the models get too complicated for this online course. you
may try our Physics_T runtime book if you need to model up to three-body collisions in
three dimensions
In two dimensions the final velocities have two components each, an x and a y velocity,
so here are four unknowns to be determined. Since momentum is a vector quantity, in
two dimensions, conservation of momentum gives us two equations, one in x and one in
y. Kinetic energy is a scalar quantity so the conservation of kinetic energy only gives us
only one equation, leaving us one equation short of the number we need to solve for four
unknowns.
What we need here is some additional information. To get that we have to look at the
geometry of the collision itself. Since we are dealing with spherically symmetrical
objects that geometry is not too hard to work out. The velocity of each ball is a vector
which can be resolved into components one of which is on the line connecting the
centers and the other perpendicular to this line. At the instant of the collision, the
component in the line of sight, so to speak, will be the only component producing a force
on the other ball. The perpendicular component can not change a ball's momentum
because it does not act through the center of mass at all (Think about it.) Likewise the
collision cannot change the velocity across the line of sight.
So in a sense all collisions are one dimensional if the dimension is carefully chosen at
the moment of impact. We can calculate the velocities in the line of sight at the instant of
impact from the initial velocities of the balls. Then we can apply the one-dimensional
analysis to determine what the velocities in the line of sight will be immediately after the
collision. Next we take the change in line of sight velocities from before to after the
collision, and recognizing that no other component of the initial velocities will be
affected by the collision, we add this change for each ball to the initial velocities to get
the after collision velocities. Run the Collision Geometry display to see the situation. A
bit later you will be able to experiment with collisions where you will control the masses
and the velocities.

Let's just briefly revisit the notion of elasticity. So far we have been assuming perfect
elasticity, meaning that kinetic energy was conserved. In real collisions, except for some
of those between nuclear particles, the elasticity is less than 100%. For now we will just
accept that there may be some loss of kinetic energy between the before collision and
after collision system. Even a 0% elastic collision does not mean that the objects come to
a stop, they just stick together. So the loss in kinetic energy can only be as much as is
consistent with the conservation of momentum. I have not been able to write a good
simple model for inelastic collisions. When I solve that problem, we may be able to go
further with this discussion.
Now we can get a bit more precise with a billiard type display. Actually we are going to
use two balls and a box in this model. Sort of like a billiard table with the possibility of
using a variety of balls in a set up shot. We will not consider the effect of rotation. My
computer is not fast enough to deal with the math involved in that. Run the Collision
Model display to try out different values of the mass, speed and directness for two
dimensional collisions.
Centroid
The centroid of a body is the center of its mass (or masses), the point at which it would
be stable, or balance, under the influence of gravity. There are other names for the same
point. It is also often called the center of gravity and the geocenter and barycenter.
There are three common "centers of gravity" that are studied in math, science and
engineering. The most common in math is the center of masses located at the vertices of
a polygon. This is more common because the other two cases can be reduced to a
variation of this approach. It is this case of point masses at the vertices that I mean when
I use centroid or center of gravity in this note, unless otherwise stated. A second
approach is to treat the area of the polygon as if it were a sheet of uiniform density. The
third, and least common, approach is to represent the sides of the polygon as wire rods of
uniform density.
Most students are first introduced to the terms above in reference to a point in a triangle.
Since the center of masses at the vertices in a triangle give the same point as a uniform
sheet, they are often confused about the various distinctions. The three centers of gravity

are usually different points in other non-symmetric polygons. It is this point, the center
of balance for the uniform sheet and also of point masses at the vertices, that is almost
universally referenced as the centroid of a triangle.
The centroid of a triangle is a point at the intersection of the three medians of the
triangle. One of the basic ideas known about the centroid is that it it divides the
medians into a 2:1 ratio. The part of the median nearest the vertex is always twice as
long as the part near the midpoint of the side. If the coordinates of the triangle are
known, then the coordinates of the centroid are the averages of the coordinates of the
vertices. If we call the three vertices A=(x1,y1); B=(x2,y2) and C=(x3,y3) then the
coordinates of the geocenter would be
.
This is extendable to the centroid of a tetrahedron in three space. If we construct the
centroids of each of the triangular faces and construct the "medial segment" from each
vertex to the centroid of the opposite face, they will also intersect in a single point, the
centroid of the tetrahedron. and the x,y,z coordinates of the centroid is the average of the
corresponding coordinates of the four vertices. The centroid of a tetrahedron divides the
medial segments into a 3:1 ratio.
If the midpoints of opposite pairs of edges of a tetrahedron are connected they will all
intersect at the centroid also. Recently while rereading Great Moments in Mathematics
before 1650 by Howard Eves I found that the extension above was called
"Commandino's theorem". It is named for Federigo Commandino and was published in
his De Centro Gravitatus Solidorum [On centers of gravity of solids] in 1565.
Commandino translated many of the classic Greek mathematical texts and a quote from
St Andrews web site says that Commandino "had the greatest influence of anybody in
ensuring that the classic Greek mathematical texts survived by publishing his editions of
them."
In a quadrilateral, the line joining the midpoints of two opposite sides is called a
bimedian. The centroid of masses located at the vertices of a quadrilateral is also the
intersection of the bimedians of a quadrilateral. Another property of the quadrilaterals
centroid is that it is also the midpoint of the segment joining the midpoints of the
diagonals

The advantage of using the point mass approach to finding centers of gravity is that the
other two common cases can be reduced to point masses of uneven weights very easily.
To find centers of gravity of uniform density sheets, one can simply divide the polygon
into non-overlapping triangles and treat the system as a set of point masses at the
centroids of these triangles with a mass equal to the area of the triangle. To find the
center of uniform rods along the perimeter of a polygon, replace each side with a point
mass equal to the length of the line located at its midpoint. The center of gravity of
uniform wire rods on the perimeter of a triangle is the Spieker point, which is the
incenter of the medial triangle. Professor Kimberling has a page showing how to find the
center of mass of any shape by a physical method.
A direct method of finding the center of gravity of a uniform density sheet in the shape
of a quadrilateral was found by F. Wittenbauer (1857-1922). If the triesectors of each
edge of the quadrilateral are found, and lines are drawn through each pair of trisectors
adjacent to a vertex, they form a parallelogram, Wittenbauer's Parallelogram. The
center of the parallelogram is the center of gravity of the uniform sheet.
Centroid Of An Area
The centroid of an area is similar to the center of mass of a body. Calculating the
centroid involves only the geometrical shape of the area. The center of gravity will equal
the centroid if the body is homogenous i.e. constant density. Integration formulas for
calculating the Centroid are:
When calculating the centroid of a complex shape. Divide the shape up into a
combination of known shapes. Then use the the following formula:

The distance from the y-axis to the centroid is Cx

The distance from the x-axis to the centroid is Cy
The coordinates of the centroid are (Cx , Cy)
The centroid location of many common shapes is known. The Properties of Areas page
includes the Centroid Location, Area, Area Moments of Inertia, Area Polar Moments of
Inertia, & Area Radius of Gyration for many common shapes.
Calculating the area and centroid of a polygon
Area
The problem of determining the area of a polygon seems at best messy but the final
formula is particularly simple. The result and sample source code (C) will be presented
here. Consider a polygon made up of line segments between N vertices (xi,yi), i=0 to N-
1. The last vertex (xN,yN) is assumed to be the same as the first, ie: the polygon is closed.
Note for polygons with holes. The holes are usually defined by ordering the vertices of
the enclosing polygon in the opposite direction to those of the holes. This algorithm still
works except that the absolute value should be taken after adding the polygon area to the
area of all the holes. That is, the holes areas will be of opposite sign to the bounding
polygon area.

The sign of the area expression above (without the absolute value) can be used to
determine the ordering of the vertices of the polygon. If the sign is positive then the
polygon vertices are ordered counter clockwise about the normal, otherwise clockwise.
To derive this solution, project lines from each vertex to some horizontal line below the
lowest part of the polygon. The enclosed region from each line segment is made up of a
triangle and rectangle. Sum these areas together noting that the areas outside the polygon
eventually cancel as the polygon loops around to the beginning.
The only restriction that will be placed on the polygon for this technique to work is that
the polygon must not be self intersecting, for example the solution will fail in the
following cases.
Centroid
The centroid is also known as the "centre of gravity" or the "center of mass". The
position of the centroid assuming the polygon to be made of a material of uniform
density is given below. As in the calculation of the area above, xN is assumed to be x0, in
other words the polygon is closed.

Centroid of a 3D shell described by 3 vertex facets
The centroid C of a 3D object made up of a collection of N triangular faces with vertices

(ai,bi,ci) is given below. Ri is the average of the vertices of the i'th face and Ai is twice the
area of the i'th face. Note the faces are assumed to be thin sheets of uniform mass, they
need not be connected or form a solid object. This reduces to the equations above for a
2D 3 vertex polygon.
Centroids & Moment of Inertia
The centroid, or center of gravity, of any object is the point within that object from
which the force of gravity appears to act. An object will remain at rest if it is balanced on
any point along a vertical line passing through its center of gravity. In terms of moments,
the center of gravity of any object is the point around which the moments of the
gravitational forces completely cancel one another.
The center of gravity of a rock (or any other three dimensional object) can be found by
hanging it from a string. The line of action of the string will always pass through the
center of gravity of the rock. The precise location of the center of gravity could be
determined if one would tie the string around the rock a number of times and note each
time the line of action of the string. Since a rock is a three dimensional object, the point
of intersection would most likely lie somewhere within the rock and out of sight.
The centroid of a two dimensional surface (such as the cross-section of a structural

shape) is a point that corresponds to the center of gravity of a very thin homogeneous
plate of the same area and shape. The planar surface (or figure) may represent an actual

area (like a tributary floor area or the cross-section of a beam) or a figurative diagram
(like a load or a bending moment diagram). It is often useful for the centroid of the area
to be determined in either case.
Symmetry can be very useful to help determine the location of the centroid of an area. If
the area (or section or body) has one line of symmetry, the centroid will lie somewhere
along the line of symmetry. This means that if it were required to balance the area (or
body or section) in a horizontal position by placing a pencil or edge underneath it, the
pencil would be best laid directly under the line of symmetry.
If a body (or area or section) has two (or more) lines of symmetry, the centroid must lie
somewhere along each of the lines. Thus, the centroid is at the point where the lines
intersect. This means that if it were required to balance the area (or body or section) in a
horizontal position by placing a nail underneath it, the point of the nail would best be
placed directly below the point where the lines of symmetry meet. This might seem
obvious, but the concept of the centroid is very important to understand both graphically
and numerically. The position of the center of gravity for some simple shapes is easily
determined by inspection. One knows that the centroid of a circle is at its center and that
of a square is at the intersection of two lines drawn connecting the midpoints of the
parallel sides. The circle has an infinite number of lines of symmetry and the square has
four. (Two were described above - what are the other two lines of symmetry?)
The centroid of a section is not always within the area or material of the section. Hollow
pipes, L shaped and some irregular shaped sections all have thir centroid located outside
of the material of the section. This is not a problem since the centroid is really only used
s a reference point from which one measures distances. The exact location of the
centroid can be determined as described above, with graphic statics, or numerically.
The centroid of any area can be found by taking moments of identifiable areas (such as
rectangles or triangles) about any axis. This is done in the same way that the center of
gravity can be found by taking moments of weights. The moment of an large area about
any axis is equal to the algebraic sum of the moments of its component areas. This is
expressed by the following equation:
Sum MAtotal = MA1 + MA2 + MA3+ ...

The moment of any area is defined as the product of the area and the perpendicular
distance from the centroid of the area to the moment axis. By means of this principle, we
may locate the centroid of any simple or composite area.
The Moment of Inertia (I) is a term used to describe the capacity of a cross-section to
resist bending. It is always considered with respect to a reference axis such as X-X or Y-
Y. It is a mathematical property of a section concerned with a surface area and how that
area is distributed about the reference axis. The reference axis is usually a centroidal
axis.
The moment of inertia is also known as the Second Moment of the Area and is
expressed mathematically as:
Ixx = Sum (A)(y2)
In which:
Ixx = the moment of inertia around the x axis

A = the area of the plane of the object
y = the distance between the centroid of the object and the x axis
The Moment of Inertia is an important value which is used to determine the state of
stress in a section, to calculate the resistance to buckling, and to determine the amount of
deflection in a beam. For example, if a designer is given a certain set of constraints on a
structural problem (i.e. loads, spans and end conditions) a "required" value of the
moment of inertia can be determined. Then, any structural element which has at least
that specific moment of inertia will be able to be utilized in the design. Another example
could be in the inverse were true: a specific element is given in a design. Then the load
bearing capacity of the element could be determined.
Let us look at two boards to intuitively determine which will deflect more and why. If
two boards with actual dimensions of 2 inches by 8 inches were laid side by side - one
on the two inch side and the other on the eight inch side, the board which is supported on
its 2" edge is considerably stiffer than that supported along its 8" edge. Both boards have
the same cross-sectional area, but the area is distributed differently about the horizontal
centroidal axis.

Calculus is ordinarily used to find the moment of inertia of an irregular section.

However, a simple formula has been derived for a rectangular section which will be the
most important section for this course.
Ixx = (1/12 ) (b)(h3)
In which the value b is always taken to be the side parallel to the reference axis and h
the height of the section. This is very important to note! If the wrong value is assumed
for the value of b, the calculations will be totally wrong.
Moment of Inertia
The importance of the distribution of the area around its centroidal axis becomes clear
when comparing the values of the moment of inertia of a number of typical beam
configurations. All of the members shown below are 2" x 6"; in cross section, equal in
length and equally loaded.
2. Hydraulics
Archimedes' principle is the law of buoyancy. It states, "any body partially or

completely submerged in a fluid is buoyed up by a force equal to the weight of the fluid
displaced by the body." The weight of an object acts downward, and the buoyant force
provided by the displaced fluid acts upward. If these two forces are equal, the object

floats. Density is defined as weight per volume. If the density of an object exceeds the
density of water, the object will sink.
Archimedes' principle, a force equal to the weight of the displaced fluid buoys up
principle that states that a body immersed in a fluid. The principle applies to both
floating and submerged bodies and to all fluids, i.e., liquids and gases. It explains not
only the buoyancy of ships and other vessels in water but also the rise of a balloon in the
air and the apparent loss of weight of objects underwater. In determining whether a given
body will float in a given fluid, both weight and volume must be considered; that is, the
relative density, or weight per unit of volume, of the body compared to the fluid
determines the buoyant force. If the body is less dense than the fluid, it will float or, in
the case of a balloon, it will rise. If the body is denser than the fluid, it will sink. Relative
density also determines the proportion of a floating body that will be submerged in a
fluid. If the body is two thirds as dense as the fluid, then two thirds of its volume will be
submerged, displacing in the process a volume of fluid whose weight is equal to the
entire weight of the body. In the case of a submerged body, the apparent weight of the
body is equal to its weight in air less the weight of an equal volume of fluid. The fluid
most often encountered in applications of Archimedes' principle is water, and the
specific gravity of a substance is a convenient measure of its relative density compared
to water. In calculating the buoyant force on a body, however, one must also take into
account the shape and position of the body. A steel rowboat placed on end into the water
will sink because the density of steel is much greater than that of water. However, in its
normal, keel-down position, the effective volume of the boat includes all the air inside it,
so that its average density is then less than that of water, and as a result it will float.
Examples: We will consider three different examples where buoyant forces play an
important role.
A nearly weightless but rigid 1-liter milk carton is held submerged under water.
What is the volume of the carton in m3? We know that 1,000 liter = 1m3. So

What is the volume of water displaced?
Since the milk carton is completely submerged, the volume of water displaces is equal to
the volume of the carton.
Volume of water displaced = 1 liter = 0.001 m3

What mass of water is displaced?
Here we use the definition of density and the known density of water (1000 kg/m 3).
What weight of water is displaced?
Now use the definition of weight.
W = mg = (1kg) (10m/sec2) = 10 N
What force is required to hold the milk carton underwater?
In this example the force that must be applied is equal and opposite to the upward
buoyant force. Since the buoyant force is simply the weight of the displaced fluid we
have F = BF = Weight of displaced water = 10 N
Buoyancy
If one places a copper ball in a pail of water it will sink, whereas a wooden ball will
float. The buoyant force on the object determines whether or not a given object will sink
or float in a fluid. The difference between the pressure at the top of the object, which
pushes it downward, and the pressure at the bottom, which pushes it upward, essentially
cause the buoyant force. Since the pressure at the bottom is always greater than at the

top, every object submerged in a fluid necessarily feels an upward buoyant force. Of
course, objects also feel a downward force due to gravity, and the difference between the
gravitational force and buoyant force on a submerged object determines whether that
object will sink, or rise to the surface. If the weight is greater than the buoyant force, the
object sinks, and vice versa. It was Archimedes (supposedly while in his bath), who
realized that submerged objects always displace fluid upwards (the level of water in the
bathtub rose when Archimedes got in). Thus, he reasoned that the buoyant force on an
object must be equal to the weight of fluid that object displaces. If the weight of an
object is greater than the weight of displaced fluid, it will sink, whereas if the weight of
the object is less than the weight of displaced fluid, it will rise. Moreover, it is evident
that the volume of displaced fluid is precisely equal to the volume of the submerged part
of the object, so that the relative density of the object and the fluid determines the
difference between the buoyant force and the weight.
To see how buoyancy works, consider a submerged brick, of height h, width w and
length l. The difference in pressure on top and bottom of the brick is ρgh, so the
difference in total force on top and bottom of the brick is simply (ρgh)(wl) = ρgV, where
V is the volume of the brick. The forces on the sides have no vertical components, so
they do not matter. The net upward force is the weight of a volume V of the fluid of
density ρ. Any body can be considered made up of brick shapes, as small as desired, so
the result applies in general. This is just the integral calculus in action, or the application
of Professor Thomson's analogy.
Consider a man in a rowboat on a lake, with a large rock in the boat. He throws the rock
into the water. What is the effect on the water level of the lake? Suppose you make a
drink of ice water with ice cubes floating in it. What happens to the water level in the
glass when the ice has melted?

The force exerted by the water on the bottom of a boat acts through the centre of gravity
B of the displaced volume, while the force exerted by gravity on the boat acts through its
own center of gravity A. This looks bad for the boat, since the boat's c.g. will naturally
be higher than the e.g. of the displaced water, so the boat will tend to capsize. Well, a
board floats, and can tell us why. Should the board start to rotate to one side, the
displaced volume immediately moves to that side, and the buoyant force tends to correct
the rotation. A floating body will be stable provided the line of action of the buoyant
force passes through a point M above the c.g. of the body, called the metacentre, so that
there is a restoring couple when the boat heels. A ship with an improperly designed hull
will not float. It is not as easy to make boats as it might appear.
Archimedes's Principle can also be applied to balloons. The Montgolfier brothers' hot air
balloon with a paper envelope ascended first in 1783 (the brothers got Pilâtre de Rozier
and Chevalier d'Arlandes to go up in it). Such "fire balloons" were then replaced with
hydrogen-filled balloons, and then with balloons filled with coal gas, which was easier to
obtain and did not diffuse through the envelope quite as rapidly. Methane would be a
good filler, with a density 0.55 that of air. Slack balloons, like most large ones, can be
contrasted with taut balloons with an elastic envelope, such as weather balloons. Slack
balloons will not be filled full on the ground, and will plump up at altitude. Balloons are
naturally stable, since the center of buoyancy is above the center of gravity in all
practical balloons. Submarines are yet another application of buoyancy, with their own
characteristic problems.
Small neoprene or natural rubber balloons have been used for meteorological
observations, with hydrogen filling. A 10g ceiling balloon was about 17" in diameter
when inflated to have a free lift of 40g. It ascended 480ft the first minute, 670ft in a
minute and a half, and 360ft per minute afterwards, to find cloud ceilings by timing, up
to 2500ft, when it subtended about 2' of arc, easily seen in binoculars. Large sounding
balloons were used to lift a radiosonde and a parachute for its recovery. An AN/AMT-2
radiosonde of the 1950's weighed 1500g, the paper parachute 100g, and the balloon
350g. The balloon was inflated to give 800g free lift, so it would rise 700-800 ft/min to
an altitude of about 50,000 ft (15 km) before it burst. This balloon was about 6 ft in

diameter when inflated at the surface, 3 ft in diameter before inflation. The information
was returned by radio telemetry, so the balloon did not have to be followed optically. Of
intermediate size was the pilot balloon, which was followed with a theodolite to
determine wind directions and speeds. At night, a pilot balloon could carry a light for
ceiling determinations.
The greatest problem with using hydrogen for lift is that it diffuses rapidly through many
substances. Weather balloons had to be launched promptly after filling, or the desired
free lift would not be obtained. Helium is a little better in this respect, but it also diffuses
rapidly. The lift obtained with helium is almost the same as with hydrogen (density 4
compared to 2, where air is 28.97). However, helium is exceedingly rare, and only its
unusual occurrence in natural gas from Kansas makes it available. Great care must be
taken when filling balloons with hydrogen to avoid sparks and the accumulation of
hydrogen in air, since hydrogen is exceedingly flammable and explosive over a wide
range of concentrations. Helium has the great advantage that it is not inflammable.
The hydrogen for filling weather balloons came from compressed gas in cylinders, from
the reaction of granulated aluminium with sodium hydroxide and water, or from the
reaction of calcium hydroxide with water. The chemical reactions are 2Al + 2NaOH +
2H2O → 2NaAlO2 + 3H2, or CaH2 + 2H2O → Ca(OH) 2 + 2H2. In the first, silicon or zinc
could be used instead of aluminium, and in the second, any similar metal hydride. Both
are rather expensive sources of hydrogen, but very convenient when only small amounts
are required. Most hydrogen is made from the catalytic decomposition of hydrocarbons,
or the reaction of hot coke with steam. Electrolysis of water is an expensive source, since
more energy is used than is recovered with the hydrogen. Any enthusiasm for a
"hydrogen economy" should be tempered by the fact that there are no hydrogen wells,
and all the hydrogen must be made with an input of energy usually greater than that
available from the hydrogen, and often with the appearance of carbon. Although about
60,000 Btu/lb is available from hydrogen, compared to 20,000 Btu/lb from gasoline,
hydrogen compressed to 1000 psi requires 140 times as much volume for the same
weight as gasoline. For the energy content of a 13-gallon gasoline tank, a 600-gallon

hydrogen tank would be required. The critical temperature of hydrogen is 32K, so liquid
storage is out of the question for general use
Practical working of: floating objects such ships

What Floats Your Boat?
Synopsis
Students use modeling clay, a material that is denser than water and thus
ordinarily sinks in water, to discover the principle of buoyancy. This exercise is intended
to follow the exercises Clay Boats and Floaters and Sinkers.
Extensions
Ask students if they themselves are floaters or sinkers. Ordinarily, humans float,
but barely. Our bodies contain mostly water, but we do have minerals in our bones,
which are denser than water, and air in our lungs (even when we exhale fully, there is
still some residual air inside), which is much less dense than water. No matter how thin
we are, we all still have some amount of fat in our bodies, which is less dense than water.
How easily we float depends largely on how much body fat we have. We can also float
better if we fill our lungs and hold our breath.
You can also ask your students if it is easier to float in a lake (fresh water) or in
the ocean (salt water). They should remember, if they have had any experience, that it is
easier to float in the latter. Ask them why this is true, and give them some time to try to
work out this puzzle. You may need to remind them what density really means; it is the
amount of stuff packed into a given space, and liquids and gases have densities, too.
Ocean water has salt, a type of "stuff", dissolved into the water. With more stuff in it,
ocean water is denser than fresh water. Put another way, a gallon of salt water weighs
more than a gallon of fresh water. Because of this weight difference, when an object
floats in salt water, a smaller volume of water needs to be displaced than would be
needed if the water contained no salt. Because less salt water needs to be displaced, an

object floats higher in the salt water. We humans find it easier to keep our heads out of
the water when we float in salt water because we don't need to displace quite as much
water as we would in fresh water.
Your students can design some simple experiments to compare the volumes of
displaced liquid when clay boats float in fresh versus salt water. (Sea water contains
about 3.5% salt, so students can make 500 mL of "artificial" sea water by combining
17.5 grams of table salt (or NaCl) with 482.5 mL of water.) You can also ask your
students what would happen when a fully loaded cargo ship travels from the ocean up
into a large river such as the Amazon. They can simulate this experience by making a
clay boat and filling it with as many metal washers (from the Clay Boats exercise) as it
will hold without sinking in salt water. Then they can carefully move the boat and its
load into a container of fresh water, where it will most likely sink.
If students have ever swum in a deep lake in the summer, they may have
discovered that the surface water was warm, but if they dove down a few feet, the water
became cooler. The colder, denser water remains at the bottom, unless something
physically causes the warm and cool water to mix. (In a relatively shallow swimming
pool, the action of swimmers is enough to keep the water mixed.) Students can use their
clay boats to try to compare the densities of very hot and very cold water. However, the
density differences are so small that students will probably not be able to measure the
slight differences in the displaced volumes accurately enough to notice any density
differences. (There is only a 1.4% difference between the density of ice water at 4° C
and hot tap water at 55° C.) That doesn't mean students shouldn't attempt the
experiment! Their results would, in fact, lead them to conclude that no such density
differences exist. They could then be asked to do some library research to find the
densities of water at different temperatures, and once that was accomplished, ask them if
their experimental set up would have allowed them to detect the reported differences.
There is a useful lesson about the limitations of measuring methods here. You can ask
your students what kinds of instruments would be needed in order to get the required
accuracy.

You can also ask students to speculate about what would happen to an oil tanker
that was loaded with crude oil from the very cold and very salty northern Atlantic Ocean,
and then traveled to the tropical, fresh water Amazon River to deliver it. (The ship
would sink lower in the water in the Amazon.) Students might be interested to know
about the Plimsoll lines painted on the hulls of ships, which indicate how fully ships can
be safely loaded under a variety of water and climate conditions. These are shown in the
figure below
the fruits of the Industrial Revolution. Manufactured goods were being produced and
exported as never before, and unfortunately, it was not uncommon for greedy ship
owners to both overload and heavily insure their boats before launching them on their
trade routes. As a result, a great many ships were lost; in one year alone 411 ships were
wrecked along England's coast, and over 500 sailors drowned. Not surprisingly, sailors
in the Merchant Navy grew reluctant to work on what they considered dangerous ships.
Many were more willing to spend twelve weeks in prison for disobeying orders than
they were to set sail on such vessels. At one time a prison inspector in southwest
England found that nine out of twelve inmates were sailors who had been jailed for
desertion.
Samuel Plimsoll, a Member of Parliament, became concerned about what he

called the "coffin ships" of the times, and fought for several years for safety regulations.

After several defeats of his proposals, Parliament finally passed the Merchant Shipping
Act of 1876, which made load lines painted on the sides of ships mandatory. The
overloading of ships was not as common in the U.S., and it wasn't until 1929 that
Congress passed the Load Line Act here. Since then, the practice of applying load
markings has become worldwide.
The basic load line symbol consists of a circle with a horizontal line passing
through the center. The letters on either side of the circle indicate the organization that
certified the load level. In the example above, the A and B stand for the American
Bureau of Shipping. The placement of the line depends on a lot of factors, including
various aspects of the hull shape, the length of the boat, and the type of cargo. The
additional lines to the right of the circular symbol show the limits for different water
types and climates, reflecting the density differences present at the time of loading.
Capillarity
Capillarity or capillary action, phenomenon in which the surface of a liquid is observed

to be elevated or depressed where it comes into contact with a solid. For example, the
surface of water in a clean drinking glass is seen to be slightly higher at the edges, where
it contacts the glass, than in the middle. Capillarity can be explained by considering the
effects of two opposing forces: adhesion, the attractive (or repulsive) force between the
molecules of the liquid and those of the container, and cohesion, the attractive force
between the molecules of the liquid (see adhesion and cohesion). Adhesion causes water
to wet a glass container and thus causes the water's surface to rise near the container's
walls. If there were no forces acting in opposition, the water would creep higher and
higher on the walls and eventually overflow the container. The forces of cohesion act to
minimize the surface area of the liquid (see surface tension); when the cohesive force
acting to reduce the surface area becomes equal to the adhesive force acting to increase it
(e.g., by pulling water up the walls of a glass), equilibrium is reached and the liquid
stops rising where it contacts the solid. In some liquid-solid systems, e.g., mercury and
glass or water and polyethylene plastic, the liquid does not wet the solid, and its surface

is depressed where it contacts the solid. Capillarity is one of the causes of the upward
flow of water in the soil and in plants.
Surface Tension
The cohesive forces between liquid molecules are responsible for the phenomenon
known as surface tension. The molecules at the surface do not have other like molecules
on all sides of them and consequently they cohere more strongly to those directly
associated with them on the surface. This forms a surface "film" which makes it more
difficult to move an object through the surface than to move it when it is completely
submersed.
Surface tension is typically measured in dynes/cm, the force in dynes required to break a
film of length 1 cm. Equivalently, it can be stated as surface energy in ergs per square
centimeter. Water at 20°C has a surface tension of 72.8 dynes/cm compared to 22.3 for
ethyl alcohol and 465 for mercury.
Cohesion and Surface Tension

The cohesive forces between molecules down into a liquid are shared with all
neighboring atoms. Those on the surface have no neighboring atoms above, and exhibit
stronger attractive forces upon their nearest neighbors on the surface. This enhancement
of the intermolecular attractive forces at the surface is called surface tension.

Surface Tension of Water

The surface tension of water is 72 dynes/cm at 25°C . It would take a force of 72 dynes
to break a surface film of water 1 cm long. The surface tension of water decreases
significantly with temperature as shown in the graph. The surface tension arises from the
polar nature of the water molecule. Hot water is a better cleaning agent because the
lower surface tension makes it a better "wetting agent" to get into pores and fissures
rather than bridging them with surface tension. Soaps and detergents further lower the
surface tension
Cohesion and Adhesion

Molecules liquid state experience strong intermolecular attractive forces. When those
forces are between like molecules, they are referred to as cohesive forces. For example,

the molecules of a water droplet are held together by cohesive forces, and the especially
strong cohesive forces at the surface constitute surface tension.
When the attractive forces are between unlike molecules, they are said to be adhesive
forces. The adhesive forces between water molecules and the walls of a glass tube are
stronger than the cohesive forces lead to an upward turning meniscus at the walls of the
vessel and contribute to capillary action.
The attractive forces between molecules in a liquid can be viewed as residual
electrostatic forces and are sometimes called van der Waals forces or van der Waals
bonds.
Water Vapor – Facts and Mechanism
a. Water vapor can be found in the air as humidity with variable degrees of saturation.
Air is called saturated when it cannot carry any more water vapor at that specific
temperature and atmospheric pressure.
b. In a saturated air at 100 ° C. at the normal atmospheric pressure, the water vapor
replaces the air totally. Any air introduced to this saturated air will change the balance
and will have different degrees of saturation. For example, at 98 ° C One cu m of air can
carry 559 grams of water vapor, and at 15 °C (the temperature of water at 14 meters
under sea level most of the year in Southern India) the vapor amount that makes One cu
m of air saturated will be 13 grams.

c. Condensation of water vapor will occur when the temperature of the condensing
body is lower than the temperature of the air at saturation. For example, at 15 ° C if the
vapor content of the air is 16 grams per cubic meter, then 3 grams will no doubt
condense.
d. The heat of vaporization will have to be released when the condensation occurs. Plus
the heat that is absorbed by the water vapor from the pounding solar energy during the
period when it leaves the water surface until it enters the pedestal (although short period)
will over a period of time, raise the temperature of the water vapor (steam) to more than
100 ° C, this additional heat will be released to the body of pipes under the cells by
conduction and not by condensation. It is true that this additional heat is marginal. But, it
is also true that the accumulation of this heat over a period of time, will make the layers
under the cells boil and add to the productivity of the plant.

FLIUD PRESSURE
Fluid pressure is the pressure on an object submerged in a fluid, such as water. The
pressure can be provided from a number of sources:
the shear weight of the fluid, such as in scuba diving, when the diver goes deeper into
the water, the water pressure increases; or in the earth's atmosphere, as a plane goes
higher, the air pressure decreases;
a pump, such as when water "pumped" into a water tower; or
a compressor, such as in a small water supply system in a rural well for a house
connected to an air compressor.
Fluid pressure occurs in one of two situations: (1) an open condition, such as the ocean,
or a swimming pool, or (2) a closed condition, such as a water line or a gas line.
Open conditions are considered to be "static" or not moving (even in the ocean where
there are waves and currents) because the fluid is essentially "at rest." The pressure in
open conditions conform with principles of fluid statics.
Closed bodies of fluid are either "static," when the fluid is not moving, or "dynamic,"
when the fluid is moving, like through a pipe. The pressure in closed conditions conform
with the principles of fluid dynamics.
The concepts of fluid pressure are predominately attributed to the discoveries of Blaise
Pascal and Daniel Bernoulli. These concepts are given in greater detail in the remainder
of this article.
Concepts of fluid pressure
There are two components to pressure: force and area. The more force, the larger the
pressure; the more area, the smaller the pressure. As such, the force on an object in
contact with fluid depends on three factors (as summarized in Bernoulli's equation):
Is the fluid at rest or moving? If moving, how fast and in what direction in relation to the
object?
What is the vertical distance from the surface of the fluid (if a liquid) to the object?
Are there any sources of pressure, such as from a pump or compressor?
Hydrostatic pressure
In the case where the fluid is at rest, called fluid statics or hydrostatics (from hydro
meaning "water" and static meaning "at rest"), the force acting on the object is the sheer
weight of the fluid above, up to the water's surface—such as from a water tower. The
resulting hydrostatic pressure is isotropic: the pressure acts in all directions equally,
according to Pascal's law:
where:
ρ (rho) is the density of the fluid (the practical density of fresh water is 1000 kg/m3);
g is the acceleration due to gravity (practical value 10 m/s2 );
h is the height of the water column in meters.
In the following example. Assume a column of 1 meter high, the pressure on the floor of
the column will be as follow:
Pa.
The pressure increases linear with the water depth.

Fluid dynamic pressure
Next, when the fluid is moving in a pipe, called fluid dynamics or hydrodynamics
(dynamics meaning "moving" or "changing"), the interaction of the fluid with the walls
of the pipe creates friction and energy loss. In fluid dynamics, the fluid pressure is
anisotropic (not isotropic), i.e. the pressure is not the same in all directions. The
anisotrophy is mainly caused by the friction losses, so the further an object is away from
the surface of the fluid (assuming the fluid is moving from the fluid surface to the
object), the less pressure there will be on that object in the fluid. (Note: There are other
reasons for the anisotrophy that are beyond the scope of this article.) These concept are
mathematically explained by Bernoulli's equation, and is applicable to all fluids.
Consider this simplified example: assume a 10 m high water tower with a large reservoir,
that would give us a static pressure of 100,000 Pa (100 kPa or 100,000 N/m2 at ground
level. Now connect a 100-meter long horizontal pipe to the bottom of the column. From
experiment it is known that the friction loss in this typical pipe, depending on material
and diameter, is 100 Pa per meter of pipe length. So for the 100 meter pipe, that would
be 10 kPa of pressure loss. So the pressure at the end of our pipe would be 100 kPa - 10
kPa = 90 kPa.
Pressure source
Finally, when a pressure source (such as an air pressure tank, a pump, or the human
heart) is attached to the system, that pressure adds to the static pressure. In very small
systems, such as in the human body or in a domestic water supply system, all of the
pressure is derived from the pressure source (the heart or an air compressor), and none
(or very little) is due to fluid static force.

Fluid Pressure Measurement
Since static fluid pressure is determined by the fluid density and depth, the depth or
height difference of a given liquid is commonly used for pressure measurement:
The fact that the liquid levels are the same in these three manometers shows that the
pressure in the glass manifold above them is uniform. This is under static conditions
with no air flow through the system so that all parts of it are at atmospheric pressure. The
fact that open liquids will seek a common level is the principle behind liquid levels for
construction purposes

The flow in this system is from left to right, driven by a high pressure air supply. The
system is said to have a positive gauge pressure exerted by the air supply. This pressure
acts to push the manometer levels down, so the minimum height in the left manometer
indicates that it is subjected to the highest pressure.
The fact that the center manometer has a higher level under these conditions of rapid air
flow indicates that the pressure has been lowered in the constriction by the Bernoulli
effect. Note that the liquid level in the right hand tube is slightly higher than the left
tube, indicating that the pressure there is slightly less than that at the left hand tube. This
is indicative of the normal pressure drop in a flowing fluid from Poiseuille's law since
the diameters of the top tube are the same at the left and right tubes

Atmospheric Pressure
The surface of the earth is at the bottom of an atmospheric sea. The standard atmospheric
pressure is measured in various units:
1 atmosphere = 760 mmHg = 29.92 inHg = 14.7 lb/in2 = 101.3 KPa
The fundamental SI unit of pressure is the Pascal (Pa), but it is a small unit so kPa is the
most common direct pressure unit for atmospheric pressure. Since the static fluid
pressure is dependent only upon density and depth, choosing a liquid of standard density
like mercury or water allows you to express the pressure in units of height or depth, e.g.,
mmHg or inches of water. The mercury barometer is the standard instrument for
atmospheric pressure measurement in weather reporting. The decrease in atmospheric
pressure with height can be predicted from the barometric formula.
The unit mmHg is often called torr, particularly in vacuum applications: 760 mmHg =
760 torr For weather applications, the standard atmospheric pressure is often called 1 bar
or 1000 millibars. This has been found to be convenient for recording the relatively
small deviations from standard atmospheric pressure with normal weather patterns.

Measurement of Pressure with the Manometer
Pressure is defined as a force per unit area - and the most accurate way to measure low
air pressure is to balance a column of liquid of known weight against it and measure the
height of the liquid column so balanced. The units of measure commonly used are inches
of mercury (in. Hg), using mercury as the fluid and inches of water (in. w.c.), using
water or oil as the fluid.
Fig. 2-1. In its simplest form the manometer is a U-tube about half filled with liquid.
With both ends of the tube open, the liquid is at the same height in each leg.
Fig. 2-2. When positive pressure is applied to one leg, the liquid is forced down in that
leg and up in the other. The difference in height, "h," which is the sum of the readings
above and below zero, indicates the pressure.
Fig. 2-3. When a vacuum is applied to one leg, the liquid rises in that leg and falls in the
other. The difference in height, "h," which is the sum of the readings above and below
zero, indicates the amount of vacuum.

Instruments employing this principle are called manometers. The simplest form is the
basic and well-known U-tube manometer. (Fig. 2-1). This device indicates the difference
between two pressures (differential pressure), or between a single pressure and
atmosphere (gage pressure), when one side is open to atmosphere. If a U-tube is filled to
the half way point with water and air pressure is exerted on one of the columns, the fluid
will be displaced. Thus one leg of water column will rise and the other falls. The
difference in height "h" which is the sum of the readings above and below the half way
point, indicates the pressure in inches of water column.
Fig. 2-4. At left, equal pressure is imposed on the fluid in the well and in the indicating
tube. Reading is zero. At the right, a positive pressure has been imposed on the liquid in
the well causing the level to go down very slightly. Liquid level in indicating tube has
risen substantially. Reading is taken directly from scale at liquid level in indicating tube.
The scale has been compensated for the drop in level in the well.
The U-tube manometer is a primary standard because the difference in height between
the two columns is always a true indication of the pressure regardless of variations in the
internal diameter of the tubing. This principle makes even the Dwyer Slack-Tube® roll-
up manometer as accurate as a laboratory instrument. This provides a real convenience
to the person who might otherwise have to board an airplane carrying a 60" long rigid
glass U-tube manometer.

VARIATIONS IN MANOMETER DESIGN

To overcome the U-tube requirement of readings at two different places, the well-type
manometer was developed. See Fig. 2-4. The reservoir (well) may be made large enough
so that the change of level in the reservoir is negligible, or the scale may be compensated
for the change in reservoir liquid level. For purposes of a more practical instrument the
Dwyer well-type manometer uses a precision bored well that requires approximately a
10% scale correction for well drop effect, thus avoiding an overly large and bulky
reservoir.
Fig. 3-1. At left, equal pressure is imposed on the liquid in the well and the indicating
tube. Reading is zero. At the right a positive pressure has been imposed on the liquid in
the indicating tube pushing it down to a point on the scale equal to the pressure. Liquid
level in the well rises proportionately. Inclining the indicating tube has opened up the
scale to permit more precise reading of the pressure.
To improve and expand readability, certain Dwyer U-type and well-type manometers are
available with a .826 sp. gr. red oil indicating fluid, and scales compensated to read
pressure directly in inches of water. To further increase readability and sensitivity the
well-type manometer indicating tube is inclined, as in Fig. 3-1, to cause a greater linear
movement along the tube for a given pressure difference. The inclined manometer is
frequently called a Draft Gage because it is widely used for determining the over-fired
draft in boiler uptakes and flues.
For an inclined manometer to be a primary device, the inclined tube must be straight and
uniform. Dwyer's precision machined solid plastic construction has been applied to a

basic line of rugged manometers, inclined and inclined-vertical, which are industry
accepted as primary instruments
Fig. 3-2. At left with equal pressure on liquid in well and indicating tube, reading is zero.
When positive pressure is imposed on liquid in indicating tube, liquid level is depressed
in tube and rises slightly in well. Reading is direct since scale is compensated for change
of level in well.
The combination of an inclined and a vertical manometer is very useful in air movement
determination. See Fig. 3-2. For air velocity measurement, an inclined scale, generally
up to 1" w.c. is used (1" w.c. velocity pressure = 4000 fpm). In the Dwyer Durablock®
inclined-vertical instrument, this scale is combined with a vertical section allowing
readings of high pressures, usually 1" w.c. to 5 to 10" w.c., to be taken. There vertical
section is used primarily for determining static pressure above the range of the inclined
section. Many special purpose types of manometers exist. Examples are the Dwyer Hook
Gage and Microtector®. These are simply U-tube manometers modified so the liquid
level can be read with a micrometer, yet retaining the basic "Physics" of the hydrostatic
U-tube primary standard. Readings accurate to ±.001" w.c. in a range of differential
pressures from 0-24" w.c. are accomplished with Dwyer Model No. 1425-24 Hook
Gage. The Model 1430 Microtector® incorporates modern electronics to increase the
accuracy of readings to ±.00025" w.c. on a 2" w.c. scale.

Center of Pressure (GRF Application Point)
As shown in Figure 1, all the forces acting between the foot and the ground can be
summed to yield a single ground reaction force vector (F) and a free torque vector (Tz).
The point of application of the ground reaction force on the plate is the center of
pressure (CP). All the small reaction forces collectively exert on the surface of the plate
at the CP.
Figure 1

Total force on immersed surfaces
Suppose we want to know the force exerted on a vertical surface of any shape with water
on one side, assuming gravity to act, and the pressure on the surface of the water zero.
We have already solved this problem by a geometrical argument, but now we apply
calculus, which is easier but not as illuminating. The force on a small area dA a distance
x below the surface of the water is dF = p dA = ρgx dA, and the moment of this force
about a point on the surface is dM = px dA = ρgx 2 dA. By integration, we can find the
total force F, and the depth at which it acts, c = M / F. If the surface is not symmetrical,
the position of the total force in the transverse direction can be obtained from the integral
of dM' = ρgxy dA, the moment about some vertical line in the plane of the surface. If
there happens to be a pressure on the free surface of the water, then the forces due to this
pressure can be evaluated separately and added to this result. We must add a force equal
to the area of the surface times the additional pressure, and a moment e qual to the
product of this force and the distance to the centroid of the surface.
The simplest case is a rectangular gate of width w, and height h, whose top is a distance
H below the surface of the water. In this case, the integrations are very easy, and F =
ρgw[(h + H)2 - h2]/2 = ρgH(H + 2h)/2 = ρg(h + H/2)Hw. The total force on the gate is
equal to its area times the pressure at its centre. M = ρgw[(h + H) 3 - h3]/3 = ρg(H2/3 + Hh
+ h2)Hw, so that c = (H 2/3 + Hh + h2)/(h + H/2). In the simple case of h = 0, c = 2H/3, or
two-thirds of the way from the top to the bottom of the gate. If we take the atmospheric
pressure to act not only on the surface of the water, but also the dry side of the gate, there

is no change to this result. This is the reason atmospheric pressure often seems to have
been neglected in solving subh problems.
Consider a curious rectangular tank, with one side vertical but the opposite side inclined
inwards or outwards. The horizontal forces exerted by the water on the two sides must be
equal and opposite, or the tank would scoot off. If the side is inclined outwards, then
there must be a downwards vertical force equal to the weight of the water above it, and
passing through the centroid of this water. If the side is inclined inwards, there must be
an upwards vertical force equal to the weight of the 'missing' water above it. In both
cases, the result is demanded by ordinary statics. What we have here has been called the
'hydrostatic paradox.' It was conceived by the celebrated Flemish engineer Simon Stevin
(1548-1620) of Brugge, the first modern scientist to investigate the statics of fluids and
solids. Consider three tanks with bottoms of equal sizes and equal heights, filled with
water. The pressures at the bottoms are equal, so the vertical force on the bottom of each
tank is the same. But suppose that one tank has vertical sides, one has sides inclined
inward, and third sides inclined outwards. The tanks do not contain the same weight of
water, yet the forces on their bottoms are equal! I am sure that you can spot the
resolution of this paradox.
Sometimes the forces are required on curved surfaces. The vertical and horizontal
components can be found by considering the equilibrium of volumes with a plane
surface equal to the projected area of the curved surface in that direction. The general
result is usually a force plus a couple, since the horizontal and vertical forces are not
necessarily in the same plane. Simple surfaces, such as cylinders, spheres and cones,
may often be easy to solve. In general, however, it is necessary to sum the forces and
moments numerically on each element of area, and only in simple cases can this be done
analytically.

If a volume of fluid is accelerated uniformly, the acceleration can be added to the

acceleration of gravity. A free surface now becomes perpendicular to the total
acceleration, and the pressure is proportional to the distance from this surface. The
same can be done for a rotating fluid, where the centrifugal acceleration is the
important quantity. The earth's atmosphere is an example. When air moves relative to
the rotating system, the Coriolis force must also be taken into account. However,
these are dynamic effects and are not strictly a part of hydrostatics.

3. Heat
Introduction to Heat
Basically, heat is energy. It manifests itself as either molecular vibration within a
substance or a s an electromagnetic radiation. Temperature on the other hand, is simply
a measure of the amount of heat present within a material
It is not possible to speak of work in a body. Similarly, it is not possible to speak of heat
in a body. Work is either done on a body or by a body. Similarly, heat can flow from a
body or to0 a body. If a body is at a constant temperature, it has both mechanic and
thermal energy due to the molecular agitation and it is not possible to separate them.
SO , in this case, we cannot talk of heat energy. It means, if the flow of heat stops the
word heat cannot be used. It is only used when there is transfer of energy between two or
more systems. Therefore, heat can be defined as energy in transit.
Temperature of a system can be defined as the property that determines whether or not
the body is in thermal equilibrium with the neighboring systems. If a number of systems
are in thermal equilibrium, this common property of the system can be represented by a
single numerical value called temperature. It means that if two systems are not in
thermal equilibrium, they are at different temperature.
Example. In mercury in glass thermometer the pressure above the mercury column is 0
and the volume if mercury measures the temperature. If a thermometer shows a constant
reading in two systems A and B separately, it will show the same reading even when A
and B are bought in contact.
Measurement of temperature of a body accurately is one of the important branches of

heat in Physics. It also becomes necessary to measure high and low temperatures. To
make this measurement possible, it is necessary to construct a suitable scale of
temperature. The scale chosen must be precise and consistent and the temperature
measured on the scale must be accurate. Accessing the temperature of a body by mere
sense of touch or comparing the degree of hotness of a body with respect to another
body does not help in measuring the temperature quantitatively and accurately.

Expansion-Contraction Coefficients
Almost without exception, solids expand upon heating and contract on cooling.
The table given below shows the change in length (expansion or contraction) when
a unit length (1 m or 1 ft, etc) of a material is warmed or cooled one-Celsius degree.
This fractional change in length for a temperature rise of 1° is called the
COEFFICIENT OF LINEAR EXPANSION. For example, 1.00 m of steel increases 2.00
E - 5 m when it is warmed 1° C. Consider a bridge 1.00 E 3 m long. For every increase
in temperature of 1° C every meter of the bridge increases by 2.00 E-5 m. This result in a
total of 2.00 E - 2 m increase in length. When the temperature increases by 10.0
degrees the increase in length will be 10 times as much. And between winter and
summer, a change of temperature of as much as 40.0° may take place, increasing the
length of this bridge by 0.800 m (1.00 E 3 m x 40.0° C x 2.00 E - 5 /° C). If the engineer
did not allow for expansion, the force exerted by this expansion (or contraction in
winter) would result in buckling of steel girders.
In metal rods or wires, the expansion or contraction that engineers are concerned about is
essentially in one dimension -- length (if you were dealing with a block of metal, you
could observe expansion or contraction in three dimensions). In a liquid, however, the
expansion or contraction in three dimensions in always important. Therefore the figures
in the table for the expansion of liquids express the COEFFICIENT OF CUBICAL
EXPANSION the change in volume for a unit volume (liter, gallon, etc.) of a liquid
changing one-Celsius degree in temperature.
You cannot easily observe the effects of temperature change on gases unless the gas is
sealed in a container. The gases of the atmosphere are free to move when expanding or
contracting. However, experimentation has shown that unlike solids and liquids, all
gases expand and contract to the same extent with a change in temperature. That is, all
gases have the same coefficient of expansion or contraction.
The behavior of substances as the temperature changes gives us another means of

identifying them. Solids melt and liquids boil; if they are pure substances, they melt or

boil at definite temperatures. But whether or not these changes occur, matter in all states
expands or contracts with a change in temperature, and the extent of this change can be
measured. Melting points, boiling points, and the amount of expansion or contraction are
three properties useful in distinguishing among substances, but some substances neither
melt nor boil. Still other properties are needed to distinguish between them.
Problems:
1. Using the tables (below), determine what change in length or volume results when:
a. An aluminum rod 1.00 m long is warmed 1.00° C:
b. A glass rod 5.00 m long is warmed 2.00° C:
c. 10.0 m of platinum wire are cooled 5.00° C:
d. 2.00 L of grain alcohol are cooled 10° C:
e. 1.00 L of benzene is warmed 5.00° C:
2. Using the tables, to what temperature must a steel rod 1.00 m long at 20.0° C
is heated to increase in length by 1.00 E - 4 m?
3. An unmarked volume of a liquid was warmed from 20 ° C to 30 ° C. This resulted

in an increase of 0.107 L. The coefficient of cubical expansion of the liquid
was 5.1 E - 4/° C. What was the original volume of the liquid?
Anomalies of Water
As water is so common place, it is often regarded as a ‘typical’ liquid. In reality water is
most atypical as a liquid, with its properties at low temperatures quite different from its
properties when hot. It has often been stated that life depends on these anomalous
properties of water. In particular, the large heat capacity, high thermal conductivity and
high water content in organisms contribute to thermal regulation and prevent local
temperature fluctuations, thus allowing us to more easily control our body temperature.
The high latent heat of evaporation gives resistance to dehydration and considerable
evaporative cooling. Water is an excellent solvent due to its polarity, high dielectric
constant and small size, particularly for polar and ionic compounds and salts. It has

unique hydration properties towards biological macromolecules (particularly proteins

and nucleic acids) that determine their three-dimensional structures, and hence their
functions, in solution. This hydration forms gels that can reversibly undergo the gel-sol
phase transitions that underlie many cellular mechanisms [Water ionizes and allows easy
proton exchange between molecules, so contributing to the richness of the ionic
interactions in biology.
At 4°C water expands on heating or cooling. This density maximum together with the
low ice density results in (i) the necessity that all of a body of fresh water (not just its
surface) is close to 4°C before any freezing can occur, (ii) the freezing of rivers, lakes
and oceans is from the top down, so permitting survival of the bottom ecology, insulating
the water from further freezing, reflecting back sunlight into space and allowing rapid
thawing, and (iii) density driven thermal convection causing seasonal mixing in deeper
temperate waters carrying life-providing oxygen into the depths. The large heat capacity
of the oceans and seas allows them to act as heat reservoirs such that sea temperatures
vary only a third as much as land temperatures and so moderate our climate (e.g. the
Gulf stream carries tropical warmth to northwestern Europe). The compressibility of
water reduces the sea level by about 40 m giving us 5% more land. Water's high surface
tension plus its expansion on freezing encourages the erosion of rocks to give soil for our
agriculture.
Notable amongst the anomalies of water are the opposite properties of hot and cold
water, with the anomalous behavior more accentuated at low temperatures where the
properties of super cooled water often diverge from those of hexagonal ice. 3 As cold
liquid water is heated it shrinks, it becomes less easy to compress, its refractive index
increases, the speed of sound within it increases, gasses become less soluble and it is
easier to heat and conducts heat better. In contrast as hot liquid water is heated it
expands, it becomes easier to compress, its refractive index reduces, the speed of sound
within it decreases, gasses become more soluble and it is harder to heat and a poorer
conductor of heat. With increasing pressure, cold water molecules move faster but hot
water molecules move slower. Hot water freezes faster than cold water and ice melts
when compressed except at high pressures when liquid water freezes when compressed.
No other material is commonly found as solid, liquid and gas.4

The anomalies
 Water has unusually high melting point.
 Water has unusually high boiling point.
 Water has unusually high critical point.
 Water has unusually high surface tension and can bounce.
 Water has unusually high viscosity.
 Water has unusually high heat of vaporization.
 Water shrinks on melting.
 Water has a high density that increases on heating (up to 3.984°C).
 The number of nearest neighbors increases on melting.
 The number of nearest neighbors increases with temperature.
 Pressure reduces its melting point (13.35 MPa gives a melting point of -1°C)
 Pressure reduces the temperature of maximum density.
 D2O and T2O differ from H2O in their physical properties much more than might
be expected from their increased mass; e.g. they have increasing temperatures of
maximum density (11.185°C and 13.4°C respectively).
 Water shows an unusually large viscosity increase but diffusion decrease as the
temperature is lowered.
 Water's viscosity decreases with pressure (at temperatures below 33°C).
 Water has unusually low compressibility.
 The compressibility drops as temperature increases down to a minimum at about
46.5°C. Below this temperature, water is easier to compress as the temperature is
lowered.
 Water has a low coefficient of expansion (thermal expansivity).
 Water's thermal expansivity reduces increasingly (becoming negative) at low
temperatures.
 The speed of sound increases with temperature (up to a maximum at 73°C).
 Water has over twice the specific heat capacity of ice or steam.
 The specific heat capacity (CP and CV) is unusually high.
 Specific heat capacity; CP has a minimum.
 NMR spin-lattice relaxation time is very small at low temperatures.

 Solutes have varying effects on properties such as density and viscosity.

 None of its solutions even approach thermodynamic ideality; even D2O in H2O is
not ideal.
 X-ray diffraction shows an unusually detailed structure.
 Supercooled water has two phases and a second critical point at about -91°C.
 Liquid water may be supercooled, in tiny droplets, down to about -70°C. It may
also be produced from glassy amorphous ice between -123°C and - 149°C and
may coexist with cubic ice up to -63°C
 Solid water exists in a wider variety of stable (and metastable) crystal and
amorphous structures than other materials.
 Hot water may freeze faster than cold water; the Mpemba effect.
 The refractive index of water has a maximum value at just below 0°C.
 The solubilities of non-polar gases in water decrease with temperature to a
minimum and then rise.
 At low temperatures, the self-diffusion of water increases as the density and
pressure increase.
 The thermal conductivity of water is high and rises to a maximum at about
130°C.
 Proton and hydroxide ion mobilities are anomalously fast in an electric field
 The heat of fusion of water with temperature exhibits a maximum at -17°C
 The electrical conductivity of water rises to a maximum at about 230°C and then
falls.
 Warm water vibrates longer than cold water.
Specific Heat Capacity and Latent Heat

Specific heat
Heat that changes the physical state of a unit mass (one kilogram) of a substance without
causing any temperature change.
The specific latent heat of fusion of a solid substance is the heat required to change one
kilogram of it from solid to liquid without any temperature change. The specific latent

heat of vaporization of a liquid substance is the heat required to change one kilogram of
it from liquid to vapour without any temperature change.
Latent heat
Latent heat, heat change associated with a change of state or phase. Latent heat, also
called heat of transformation, is the heat given up or absorbed by a unit mass of a
substance as it changes from a solid to a liquid, from a liquid to a gas, or the reverse of
either of these changes. It is called latent because it is not associated with a change in
temperature. Each substance has a characteristic heat of fusion, associated with the solid-
liquid transition, and a characteristic heat of vaporization, associated with the liquid-gas
transition. The latent heat of fusion for ice is 80 calories per gram (see calorie). This
amount of heat is absorbed by each gram of ice in melting or is given up by each gram of
water in freezing. The latent heat of vaporization of steam is 540 calories per gram,
absorbed during vaporization or given up during condensation. For a substance going
directly from the solid to the gas state, or the reverse, the heat absorbed or given up is
known as the latent heat of sublimation.
Heat
Specific heat
A measure of the efficiency with which a substance can store this heat energy is known
as specific heat capacity, or simply the specific heat, . The greater the material's
specific heat, the more energy must be added to change its temperature. As an example,
the specific heat of water is given as , which means that 1.00

calorie of heat is necessary to raise one gram of water one degree Celsius, or 4190 joules
of heat are necessary to raise one kilogram of water one Kelvin. Specific heat values for
various other materials are listed in Table 1 below.
The heat necessary to cause a temperature change, , in a material of mass, , is given
by the equation (1)
Where is the specific heat capacity of that material, as usual, the change in temperature
of the substance is the difference between its final and initial temperatures, .

Table 1
Specific and Latent Heat Values
Specific Heat Latent Heat of Fusion
Material (cal/g °C) (J/kg K) (cal/g) (J/kg)
Aluminum 0.215 900 94.5 3.96x105
Copper 0.092 385 49.0 2.05x105
Iron 0.107 448 63.7 2.67x105
Lead 0.031 130 5.5 0.23x105
Brass 0.092 385 Unknown Unknown
Magnesium 0.245 1030 88.0 3.7x105
Zinc 0.093 390 27.0 1.1x105
Styrofoam 0.27 1131 Unknown Unknown
Air 0.240 1006 N/A N/A
Water 1.000 4190 N/A N/A
Ice 0.500 2095 79.7 3.34x105
When two bodies having different temperatures come into contact with each other, heat
energy is transferred between the bodies. Take, for example, placing a piece of hot metal
into a container of cool water. From experience we know over time that the metal sample
will become cooler, while the water and its container will become warmer, until
equilibrium temperature is reached. Put another way, according to the law of energy
conservation the total heat energy lost by the metal is the total heat energy gained by
the water and container:
(2)
The negative sign is used to maintain our sign convention for heat; we set the hot side to
be negative because energy is leaving the hot sample. Equation 2 holds true if no heat is

exchanged with the surrounding environment and if none of the materials undergo a
phase change.
To isolate this experiment from the environment we will utilize two types of
calorimeters: a simple styrofoam cup (Figure 3 below) and a more traditional wooden
calorimeter with an aluminum cup insert (Figure 4 below). You will determine the
effectiveness of both of these calorimeters in this experiment.
If the materials undergo a phase change, as shown in the above images, we can not use
Equation 2. If a phase change is experienced, for example by melting or boiling, the
internal energy of the material changes, but its temperature does not. Because this energy
change does not alter the temperature of the material, we refer to this heat as latent, or
"hidden", heat . In the case where a solid changes into a liquid (melting) we refer to
the latent heat of fusion, , which is defined as the amount of heat necessary to change
one gram of a solid to a liquid without a temperature change. Note that latent heat of
fusion literally means "hidden heat which causes melting". The heat required to
completely melt a given mass, , of a substance is given as
(3)
In Objective 3 of this experiment, you are to determine the latent heat of fusion of ice,
here; you will add some ice cubes to a water bath at room temperature. Heat is removed
from the water bath to both melt the ice and warm the melted ice water to an
equilibrium temperature. Again, from energy conservation arguments we can write

(4)
Note that while in each experiment, the initial temperatures of each material
isdifferent, they end up at the same the final temperature.
Teaching Suggestions, Activities and Demonstrations

1. Perform an activity to determine the specific heat capacity of a substance. See the
next activity for an example of how this could be done.
2. Place a measured volume of cold water into a calorimeter. Measure its initial
temperature. Heat a metal object of known mass in a beaker of boiling water long
enough so that the metal object reaches a temperature of the boiling point of
water. (When the metal object is first placed into boiling water, the water may
stop boiling. Wait until the water has reached the boiling point once again.)
Transfer the hot metal object into the calorimeter. Record the final temperature of
the water in the calorimeter after heat transfer has taken place. Determine the
specific heat capacity of the metal. Repeat using different types of metal. (Before
transferring the metal to the calorimeter, attach thread to it. Transferring the
metal using cold metal tongs acts as a heat sink.) Based on the experimental
results, qualitatively develop the Principle of Heat Exchange, for section C on
thermodynamics.
3. Perform an activity to investigate the specific latent heat of fusion (or the specific
latent heat of vaporization, or both) for a given substance.
Repeat the above experiment with ice to find the specific latent heat of ice.
4. Give students some examples of the unique physical characteristics of water and
how they have made it possible for life to exist on this planet. For example,
brainstorm to establish what kinds of things might happen if water in a large lake
froze solid in the winter, or if water froze at lower temperatures.
5. During a short break between classes, suppose a person wants to drink a cup of
coffee. Design an experiment to determine if it is better to add milk (or cream)

and then wait five minutes before drinking it, or to allow the black coffee to cool
for five minutes and then add the milk (or cream). In which situation would the
coffee be the coldest after the five minute interval?
Other variables, such as the type of container the coffee is poured into, or the
initial temperature of the coffee, and the type of coffee used (percolated, drip, or
instant), etc., should all be fully controlled.
6. Using three different metal samples having the same mass, heat them all in water
until each one reaches the boiling point of water. Carefully remove each one and
place it on a sheet of paraffin wax. Allow them all to cool. Observe how far each
sample has been imbedded into the wax. The metal with the highest specific heat
capacity will have been able to melt away the most wax and imbed itself the
deepest into the wax. This is a very simple but effective demonstration or student
group activity which reinforces this concept.
7. Using brainstorming, suggest some of the ways in which the unique physical
characteristics of water lead to a number of interesting and practical applications.
8. Waste heat from industrial processes sometimes causes thermal pollution.
Research ways in which waste heat is being put to better use than by releasing it
into the environment. Use webbing to link this activity with Science 10 or
Biology 20.
Design a passive solar water heater. Design it to optimize the amount of heat
energy that can be absorbed by a given quantity of water. Compare the efficiency
of the designs developed by different groups.

Boiling Point
When a liquid is heated, it eventually reaches a temperature at which the vapor pressure
is large enough that bubbles form inside the body of the liquid. This temperature is
called the boiling point. Once the liquid starts to boil, the temperature remains constant
until all of the liquid has been converted to a gas. The normal boiling point of water is
100oC. But if you try to cook an egg in boiling water while camping in the Rocky
Mountains at an elevation of 10,000 feet, you will find that it takes longer for the egg to
cook because water boils at only 90oC at this elevation.
In theory, you shouldn't be able to heat a liquid to temperatures above its normal boiling
point. Before microwave ovens became popular, however, pressure cookers were used to
decrease the amount of time it took to cook food. In a typical pressure cooker, water can
remain a liquid at temperatures as high as 120oC, and food cooks in as little as one-third
the normal time.
To explain why water boils at 90oC in the mountains and 120oC in a pressure cooker,
even though the normal boiling point of water is 100 oC, we have to understand why a
liquid boils. By definition, a liquid boils when the vapor pressure of the gas escaping
from the liquid is equal to the pressure exerted on the liquid by its surroundings, as
shown in the figure below.
Liquids boil when their vapor pressure is equal to the pressure exerted on the liquid by
its surroundings. The normal boiling point of water is 100 oC because this is the
temperature at which the vapor pressure of water is 760 mmHg, or 1 atm. Under normal
conditions, when the pressure of the atmosphere is approximately 760 mmHg, water
boils at 100oC. At 10,000 feet above sea level, the pressure of the atmosphere is only 526
mmHg. At these elevations, water boils when its vapor pressure is 526 mmHg, which
occurs at a temperature of 90oC.

Pressure cookers are equipped with a valve that lets gas escape when the pressure inside
the pot exceeds some fixed value. This valve is often set at 15 psi, which means that the
water vapor inside the pot must reach a pressure of 2 atm before it can escape. Because
water doesn't reach a vapor pressure of 2 atm until the temperature is 120 oC, it boils in
this container at 120oC.
Liquids often boil in an uneven fashion, or bump. They tend to bump when there aren't
any scratches on the walls of the container where bubbles can form. Bumping is easily
prevented by adding a few boiling chips to the liquid, which provide a rough surface
upon which bubbles can form. When boiling chips are used, essentially all of the bubbles
that rise through the solution form on the surface of these chips The boiling point is
defined as the temperature at which the saturated vapor pressure of a liquid is equal to
the surrounding atmospheric pressure. For water, the vapor pressure reaches the standard
sea level atmospheric pressure of 760 mmHg at 100°C. Since the vapor pressure
increases with temperature, it follows that for pressure greater than 760 mmHg (e.g., in a
pressure cooker), the boiling point is above 100°C and for pressure less than 760 mmHg
(e.g., at altitudes above sea level), the boiling point will be lower than 100°C. As long as
a vessel of water is boiling at 760 mmHg, it will remain at 100°C until the phase change
is complete. Rapidly boiling water is not at a higher temperature than slowly boiling
water. The stability of the boiling point makes it a convenient calibration temperature for
temperature scales.
At the boiling point, saturated vapor pressure equals atmospheric pressure

Boiling Point Variation Near 100 C
Values were taken from the saturated vapor pressure table for water near 100 degrees
Celsius. An empirical fit to these data values was made, and the formula obtained is
shown on the diagram. It could be considered to be reasonably valid only for a few
degrees above and below 100 °C since the curve is very non-linear.
Melting Point and Freezing Point
Pure, crystalline solids have a characteristic melting point, the temperature at which the
solid melts to become a liquid. The transition between the solid and the liquid is so sharp
for small samples of a pure substance that melting points can be measured to 0.1 oC. The
melting point of solid oxygen, for example, is -218.4oC.
Liquids have a characteristic temperature at which they turn into solids, known as their
freezing point. In theory, the melting point of a solid should be the same as the freezing

point of the liquid. In practice, small differences between these quantities can be
observed.
It is difficult, if not impossible, to heat a solid above its melting point because the heat
that enters the solid at its melting point is used to convert the solid into a liquid. It is
possible, however, to cool some liquids to temperatures below their freezing points
without forming a solid. When this is done, the liquid is said to be supercooled.
An example of a supercooled liquid can be made by heating solid sodium acetate

trihydrate (NaCH3CO2 3 H2O). When this solid melts, the sodium acetate dissolves in the
water that was trapped in the crystal to form a solution. When the solution cools to room
temperature, it should solidify. But it often doesn't. If a small crystal of sodium acetate
trihydrate is added to the liquid, however, the contents of the flask solidify within
seconds.
A liquid can become supercooled because the particles in a solid are packed in a regular
structure that is characteristic of that particular substance. Some of these solids form
very easily; others do not. Some need a particle of dust, or a seed crystal, to act as a site
on which the crystal can grow. In order to form crystals of sodium acetate trihydrate, Na +
ions, CH3CO2- ions, and water molecules must come together in the proper orientation. It
is difficult for these particles to organize themselves, but a seed crystal can provide the
framework on which the proper arrangement of ions and water molecules can grow.
Because it is difficult to heat solids to temperatures above their melting points, and
because pure solids tend to melt over a very small temperature range, melting points are
often used to help identify compounds. We can distinguish between the three sugars
known as glucose (MP = 150oC), fructose (MP = 103-105oC), and sucrose (MP = 185-
186oC), for example, by determining the melting point of a small sample.
Measurements of the melting point of a solid can also provide information about the
purity of the substance. Pure, crystalline solids melt over a very narrow range of
temperatures, whereas mixtures melt over a broad temperature range. Mixtures also tend
to melt at temperatures below the melting points of the pure solids.

1 High melting point (0°C, cf. CHCl3 -63°C)
(H2O) showing the melting (M.Pt.), boiling (B.Pt.) and triple

(T.Pt) points.a
The melting point of water is over 100 K higher than expected by extrapolation of the
melting points of other Group 6A hydrides, here shown compared with Group 4A
hydrides. It is also much higher than O2 (-219°C) or H2 (-269°C).
In ice (Ih), all water molecules participate in four hydrogen bonds (two as donor and
two as acceptor) and are held relatively static. In liquid water, some of the weaker

hydrogen bonds must be broken to allow the molecules move around. The large energy
required for breaking these bonds must be supplied during the melting process and only
a relatively minor amount of energy is reclaimed from the change in volume. The free
energy change (G=H-TS) must be zero at the melting point. As temperature is
increased, the amount of hydrogen bonding in liquid water decreases and its entropy
increases. Melting will only occur when there is sufficient entropy change to provide the
energy required for the bond breaking. The low entropy (high organization) of liquid
water causes this melting point to be high.fusion and differs for different substances. Ice,
for example, requires approximately 80 calories of heat to change each gram to water at
its melting point. Because its heat of fusion is relatively high, ice is used in refrigeration.
In freezing (the reverse process, i.e., the change from liquid to solid), heat is given off by
the substance undergoing the change, and the amount given off is the same as that
absorbed in melting
Factors affecting Boiling:

A liquid boils at a temperature at which its vapor pressure is equal to the pressure of
the gas above it. The lower the pressure of a gas above a liquid, the lower the
temperature at which the liquid will boil.
The Macroscopic View
 As a liquid is heated, its vapor pressure increases until the vapor pressure equals
the pressure of the gas above it.
 Bubbles of vaporized liquid (i.e., gas) form within the bulk liquid and then rise to
the surface where they burst and release the gas. (At the boiling temperature the
vapor inside a bubble has enough pressure to keep the bubble from collapsing.)
 In order to form vapor, the molecules of the liquid must overcome the forces of
attraction between them.
 The temperature of a boiling liquid remains constant, even when more heat is
added.
The boiling point of a liquid is the temperature at which its vapor pressure is equal
to the pressure of the gas above it. The normal boiling point of a liquid is the
temperature at which its vapor pressure is equal to one atmosphere (760 torr).
The Microscopic View
Microscopic view inside a bubble in boiling water. The diagram shows the right-
hand inner surface of the bubble. Note that water gas and liquid are in equilibrium

Factors That Affect the Boiling Point

 Pressure: when the external pressure is:
o Less than one atmosphere, the boiling point of the liquid is lower than its
normal boiling point.
o Equal to one atmosphere, the boiling point of a liquid is called the
normal boiling point.
o Greater than one atmosphere, the boiling point of the liquid is greater than
its normal boiling point.
Types of Molecules: the types of molecules that make up a liquid determine its
boiling point. If the intermolecular forces between molecules are:
o Relatively strong, the boiling point will be relatively high. Relatively

weak, the boiling point will be relatively low
Transference of heat: Convection. Convection. Radiation
Convection. This is the name given to the transference of heat through gas or
liquid media by the circulation of heated particles. For example: in nature there is
rapid upward movement of air particles that have been heated by contact with
surfaces on the Earth warmed by the sun. Electricity does not travel by
convection.
Conduction. Heat moves (is transferred) fairly slowly by this method which
involves heat energy being transferred molecule by molecule when a warm
substance is in direct contact with a colder substance. Conduction is a typical
method of transference of electrical energy and is rapid.
Radiation. This is a very rapid movement (transfer) of energy in the form of
electromagnetic waves (radio waves, heat, light, X-rays, cosmic rays), and sound
waves.

Heat is a form of energy. We use it for a lot of things, like warming our homes and
cooking our food
Heat energy moves in three ways:
1. Conduction
2. Convection
3. Radiation
Conduction occurs when energy is passed directly from one item to another. If you
stirred a pan of soup on the stove with a metal spoon, the spoon will heat up. The heat is
being conducted from the hot area of the soup to the colder area of spoon. Metals are
excellent conductors of heat energy. Wood or plastics are not. These "bad" conductors
are called insulators. That's why a pan is usually made of metal while the handle is made
of a strong plastic.

Convection is the movement of gases or liquids from a cooler spot to a warmer spot. If
a soup pan is made of glass, we could see the movement of convection currents in the
pan. The warmer soup moves up from the heated area at the bottom of the pan to the top
where it is cooler. The cooler soup then moves to take the warmer soup's place. The
movement is in a circular pattern within the pan (see picture
above). The wind we feel outside is often the result of convection
currents. You can understand this by the winds you feel near an
ocean. Warm air is lighter than cold air and so it rises. During the
daytime, cool air over water moves to replace the air rising up as
the land warms the air over it. During the nighttime, the
directions change -- the surface of the water is sometimes
warmer and the land is cooler
Radiation is the final form of movement of heat energy. The sun's light and heat cannot
reach us by conduction or convection because space is almost completely empty. There
is nothing to transfer the energy from the sun to the earth.The sun's rays travel in straight
lines called heat rays. When it moves that way, it is called radiation. When sunlight hits
the earth, its radiation is absorbed or reflected. Darker surfaces absorb more of the
radiation and lighter surfaces reflect the radiation. So you would be cooler if you wear
light or white clothes in the summer

Vapour pressure
Pressure of a vapour given off by (evaporated from) a liquid or solid, caused by atoms or
molecules continuously escaping from its surface. In an enclosed space, a maximum
value is reached when the number of particles leaving the surface is in equilibrium with
those returning to it; this is known as the saturated vapour pressure or equilibrium
vapour pressure
Water vapour is one of several gaseous constituents of air, the other principal ones being
nitrogen, oxygen and carbon dioxide. Each exerts its own partial pressure in proportion
to the amount of gas present, the sum of the pressures making up the total or barometric
pressure of the air. When there is a difference in concentration of one of these gases
between two points, there will be a corresponding difference in partial pressure. This will
cause a flow of that particular gas from the point of higher concentration to the lower.
When a partial pressure difference exists between two sides of a material, the gas
involved will diffuse through the material until the partial pressures of that gas are
equalized. The rate of diffusion will be determined by the partial pressure difference, the
length of the flow path, and the permeability to the particular gas involved of the
medium through which flow is taking place. When water changes state from a liquid to a
gas, as it does when it evaporates into the air, the water molecules in the vapour expand.
Just as air pressure is directly related to the number of gas molecules per cubic meter of
space, so vapour pressure results from the number of water vapour molecules per cubic
meter. The greater the moisture vapour contents of air, the greater the vapour pressure.
Thus vapour pressure is linearly related to absolute humidity and runs along the same
axis of the Psychrometric Chart.

The origin of saturated vapour pressure
The temperature governs the average energy of the particles in a liquid. The higher the
temperature, the higher the average energy. But within that average, some particles have
energies higher than the average, and others have energies lower than the average.
Some of the more energetic particles on the surface of the liquid can be moving fast
enough to escape from the attractive forces holding the liquid together. They evaporate.
The diagram shows a small region of a liquid near its surface
Notice that evaporation only takes place on the surface of the liquid. That's quite
different from boiling, which happens when there is enough energy to disrupt the
attractive forces throughout the liquid. That's why, if you look at boiling water, you see
bubbles of gas being formed all the way through the liquid.
If you look at water, which is just evaporating in the sun, you don't see any bubbles.
Water molecules are simply breaking away from the surface layer.
Eventually, the water will all evaporate in this way. The energy, which is lost as the
particles evaporate, is replaced from the surroundings. As the molecules in the water
jostle with each other, new molecules will gain enough energy to escape from the
surface.

The evaporation of a liquid in a closed container

Now imagine what happens if the liquid is in a closed container. Common sense tells
you that water in a sealed bottle doesn't seem to evaporate - or at least, it doesn't
disappear over time.
But there is constant evaporation from the surface. Particles continue to break away from
the surface of the liquid - but this time they are trapped in the space above the liquid.
As the gaseous particles bounce around, some of them will hit the surface of the liquid
again, and be trapped there. There will rapidly be an equilibrium set up in which the
number of particles leaving the surface is exactly balanced by the number rejoining it.

In this equilibrium, there will be a fixed number of the gaseous particles in the space
above the liquid.
When these particles hit the walls of the container, they exert a pressure. This pressure is
called the saturated vapour pressure (also known as saturation vapour pressure) of the
liquid.
Measuring the saturated vapour pressure
It isn't difficult to show the existence of this saturated vapour pressure (and to measure
it) using a simple piece of apparatus.
If you have a mercury barometer tube in a trough of mercury, at 1 atmosphere pressure

the column will be 760 mm tall. 1 atmosphere is sometimes quoted as 760 mmHg
("millimetres of mercury").
If you squirt a few drops of liquid into the tube, it will rise to form a thin layer floating
on top of the mercury. Some of the liquid will evaporate and you will get the equilibrium
we've just been talking about - provided there is still some liquid on top of the mercury.
It is only an equilibrium if both liquid and vapour are present.

The saturated vapour pressure of the liquid will force the mercury level down a bit. You
can measure the drop - and this gives a value for the saturated vapour pressure of the
liquid at this temperature. In this case, the mercury has been forced down by a distance
of 760 - 630 mm. The saturated vapour pressure of this liquid at the temperature of the
experiment is 130 mmHg.
You could convert this into proper SI units (pascals) if you wanted to. 760 mmHg is
equivalent to 101325 Pa.
A value of 130 mmHg is quite a high vapour pressure if we are talking about room
temperature. Water's saturated vapour pressure is about 20 mmHg at this temperature. A
high vapour pressure means that the liquid must be volatile - molecules escape from its
surface relatively easily, and aren't very good at sticking back on again either.
That will result in larger numbers of them in the gas state once equilibrium is reached.
The liquid in the example must have significantly weaker intermolecular forces than
water. down by a distance of 760 - 630 mm. The saturated vapour pressure of this liquid
at the temperature of the experiment is 130 mmHg.

You could convert this into proper SI units (pascals) if you wanted to. 760 mmHg is
equivalent to 101325 Pa.
A value of 130 mmHg is quite a high vapour pressure if we are talking about room
temperature. Water's saturated vapour pressure is about 20 mmHg at this temperature. A
high vapour pressure means that the liquid must be volatile - molecules escape from its
surface relatively easily, and aren't very good at sticking back on again either.
That will result in larger numbers of them in the gas state once equilibrium is reached.
The liquid in the example must have significantly weaker intermolecular forces than
water. When the space above the liquid is saturated with vapour particles, you have this
equilibrium occurring on the surface of the liquid:
The forward change (liquid to vapour) is endothermic. It needs heat to convert the liquid
into the vapour.
According to Le Chatelier, increasing the temperature of a system in a dynamic

equilibrium favours the endothermic change. That means that increasing the temperature
increases the amount of vapour present, and so increases the saturated vapour pressure

The effect of temperature on the saturated vapour pressure of water
The graph shows how the saturated vapour pressure (svp) of water varies from 0°C to
100 °C. The pressure scale (the vertical one) is measured in kilopascals (kPa). 1
atmosphere pressure is 101.325 kPa.
Saturated vapour pressure and boiling point
A liquid boils when its saturated vapour pressure becomes equal to the external pressure
on the liquid. When that happens, it enables bubbles of vapour to form throughout the
liquid - those are the bubbles you see when a liquid boils.
If the external pressure is higher than the saturated vapour pressure, these bubbles are
prevented from forming, and you just get evaporation at the surface of the liquid.
If the liquid is in an open container and exposed to normal atmospheric pressure, the
liquid boils when its saturated vapour pressure becomes equal to 1 atmosphere (or
101325 Pa or 101.325 kPa or 760 mmHg). This happens with water when the
temperature reaches 100°C.
But at different pressures, water will boil at different temperatures. For example, at the
top of Mount Everest the pressure is so low that water will boil at about 70°C.

Depressions from the Atlantic can easily lower the atmospheric pressure in the UK
enough so that water will boil at 99°C - even lower with very deep depressions.
Whenever we just talk about "the boiling point" of a liquid, we always assume that it is
being measured at exactly 1 atmosphere pressure. In practice, of course, that is rarely
exactly true.
Saturated vapour pressure and solids
Sublimation
Solids can also lose particles from their surface to form a vapour, except that in this case
we call the effect sublimation rather than evaporation. Sublimation is the direct change
from solid to vapour (or vice versa) without going through the liquid stage.
In most cases, at ordinary temperatures, the saturated vapour pressures of solids range
from low to very, very, very low. The forces of attraction in many solids are too high to
allow much loss of particles from the surface.
However, there are some which do easily form vapours. For example, naphthalene (used
in old-fashioned "moth balls" to deter clothes moths) has quite a strong smell. Molecules
must be breaking away from the surface as a vapour, because otherwise you wouldn't be
able to smell it.
Another fairly common example (discussed in detail on another page) is solid carbon
dioxide - "dry ice". This never forms a liquid at atmospheric pressure and always
converts directly from solid to vapour. That's why it is known as dry ice.
The process of evaporation in a closed container will proceed until there are as many
molecules returning to the liquid as there are escaping. At this point the vapor is said to
be saturated, and the pressure of that vapor (usually expressed in mmHg) is called the
saturated vapor pressure

4. Oscillations
Simple Harmonic Motion

Simple Harmonic Motion arises when we consider the motion of a particle whose
acceleration points towards a fixed point O and is proportional to the distance of the
particle from O (so the acceleration increases as the distance from the fixed point
increases).
As the particle moves away from the fixed point O, since the acceleration is pointing
towards O, the particle will slow down and eventually stop (at Q), before returning to O.
It will keep going and then again slow down as it reaches P before stopping at P and
returning to O once more.
The particle will therefore move between two fixed points (P and Q). The amplitude of
the motion is the distance from O to either P or Q (the distances are the same).
A particle, which moves under simple harmonic motion, will have the equation
= - 2x
Where  is a constant (note that this just says that the acceleration of the particle is
proportional to the distance from O).
You may be asked to prove that a particle moves with simple harmonic motion. If so,
you simply must show that the particle satisfies the above equation.
Further Equations
We can solve this differential equation to deduce that:
 v2 = 2(a2 - x2)
Where v is the velocity of the particle, a is the amplitude and x is the distance from O.
From this equation, we can see that the velocity is maximized when x = 0, since v 2 =
2a2 - 2x2
Hence the maximum velocity is a (put x = 0 in the above equation and take the square
root).
The period of the motion is the time it takes for the particle to perform one complete
cycle. It can be calculated using:

 T = 2 / 
If the particle is at 0 when t = 0, then the following equation also holds:
 x = asin t
If the particle is at P or Q when t = 0, then the following equation also holds:
x = acos t
Simple harmonic motion is typified by the motion of a mass on a spring when it is

subject to the linear elastic restoring force given by Hooke's Law. The motion is
sinusoidal in time and demonstrates a single resonant frequency.
The motion equation for simple harmonic motion contains a complete description of the
motion, and other parameters of the motion can be calculated from it.
The velocity and acceleration are given by
The total energy for an undamped oscillator is the sum of its kinetic energy and potential
energy, which is constant at

Simple Harmonic Motion Calculation

The motion equations for simple harmonic motion provide for calculating any parameter
of the motion if the others are known
If the period is T = s
then the frequency is f = Hz and the angular frequency = rad/s.

The motion is described by
Displacement = Amplitude x sin (angular frequency x time)
y =A x sin (
x t )
m= m x sin ( rad/s
x s )
Any of the parameters in the motion equation can be calculated by clicking on the active
word in the motion relationship above. Default values will be entered for any missing
data, but those values may be changed and the calculation repeated. The angular
frequency calculation assumes that the motion is in its first period and therefore
calculates the smallest value of angular frequency, which will match the other

parameters. The time calculation calculates the first time the motion reaches the
specified displacement, i.e., the time during the first period.
Simple Harmonic Motion: Simple Pendulum

Consider a particle subject to a linear restoring force. This could be a mass attached to a
spring, or a simple pendulum. We will be using a simple pendulum in this lab
experiment: Figure 1: The Simple Pendulum
If we consider only motions of the pendulum in one vertical plane, the moment of inertia
is: I = ml2 (1)
and the torque can be written as:
= I  = ml2 d2 d t2 = -mgl sin 
(2)
The torque is taken as negative since it is acting in a direction such as to decrease the
angle  . Then, the equation of motion can be written as:
d2 d t2 = - (g/l) sin 
(3)
If we were using a spring we would have a differential equation in terms of the linear
displacement x, but since this is a swinging system, must be used. Both types of systems
are solved in the same manner. We make an approximation to be valid in the case of
small angles (<< /2) so that sin is approximately  (in radians). Note that
theoretically, the frequency is independent of the amplitude (but remember that the
amplitude must be somewhat small since we made this assumption to solve the problem
to begin with). Since is measured in radians (angular measure per unit time) we can
express the period, or time for oscillation as:

T = 2  (l/g)
(4)
Note that the period depends on the length of the pendulum and not the initial
displacement of the bob. Because of this, pendulums have for ages been used to regulate
time in many clocks.
Forced Oscillations and Resonance

One can apply an additional force to the mass on a spring. Suppose this additional force
(AKA the driving force) has an angular frequency d .
Newton’s second law, that the mass undergoes very large amplitude vibrations if it so
happens the driving frequency matches the natural frequency of the SHM, can show it,
d = 
This phenomenon is known as resonance. Examples include: breaking a wine glass with
sound, the Tacoma Narrows bridge catastrophe, soldiers required to break rank when
crossing a bridge, etc.
The top of a spring pendulum (red circle) is moved to and fro - for example by hand; this
motion is assumed as harmonic, which means that it is possible to describe the motion
by a cosine function. The oscillations of the spring pendulum caused in this way are
called forced oscillations.
The "Reset" button brings the spring pendulum to its initial position. You can start or
stop and continue the simulation with the other button. If you choose the option
"Slow motion", the movement will be five times slower. The spring constant, the mass,
the constant of attenuation and the angular frequency of the exciting oscillation can be
changed within certain limits. In addition you can select one of three diagrams by using
the appropriate radio buttons:
The elongations of exciter and resonator as functions of time The amplitude of the
resonator's oscillation dependent on the exciter's angular frequency The phase difference
between the oscillations of exciter and resonator dependent on the exciter's angular
frequency On the whole you can see three different types of behavior:

If the exciter's frequency is very small (this means that the top of the spring pendulum is
moved very slowly), the pendulum will oscillate nearly synchronously with the exciter
and nearly with the same amplitude. If the exciter's frequency agrees with the
characteristic frequency of the spring pendulum, the oscillations of the pendulum will
build up more and more (resonance); in this case the oscillations are delayed about one
fourth of the oscillation period compared with the exciter.
If the exciter's frequency is very high, the resonator will oscillate only with a very small
amplitude and nearly the opposite phase.
If the constant of attenuation (the friction) is very small, the transient states will be
important too; therefore you have to wait some time in this case to notice the mentioned
types of behavior.

5. Sound
Introduction To Sound
Everyday your world is filled with a multitude of sounds. Sound can let you
communicate with others or let ot`hers communicate with you. It can be a warning of
danger or simply an enjoyable experience. Some sounds can be heard by dogs or other
animals but cannot be heard by humans.
Transverse Waves
The particles in the medium oscillate in a direction perpendicular to the way the wave is
traveling. A good example is a wave on a string. For transverse waves the displacement
of the medium is perpendicular to the direction of propagation of the wave. A ripple on a
pond and a wave on a string are easily visualized transverse waves
Transverse waves cannot propagate in a gas or a liquid because there is no mechanism

for driving motion perpendicular to the propagation of the wave.
Longitudinal Waves
The particles in the medium oscillate along the same direction as the way the wave is
traveling. Sound waves are longitudinal waves.
In longitudinal waves the displacement of the medium is parallel to the propagation of
the wave. A wave in a "slinky" is a good visualization. Sound waves in air are
longitudinal waves.

Transverse Waves & Longitudinal Waves

Most kinds of waves are transverse waves. In a transverse wave, as the wave is moving
in one direction, it is creating a disturbance in a different direction. The most familiar
example of this is waves on the surface of water. As the wave travels in one direction -
say south - it is creating an up-and-down (not north-and-south) motion on the water's
surface. This kind of wave is very easy to draw; a line going from left-to-right has up-
and-down wiggles. So most diagrams of waves - even of sound waves - are pictures of
transverse waves.
But sound waves are not transverse. Sound waves are longitudinal waves. If sound
waves are moving south, the disturbance that they are creating is making the air
molecules vibrate north-and-south (not east-and-west, or up-and-down. This is very
difficult to show clearly in a diagram, so most diagrams, even diagrams of sound waves,
show transverse waves.
Longitudinal waves may also be a little difficult to imagine, because there aren't any
examples that we can see in everyday life. A mathematical description might be that in
longitudinal waves, the waves (the disturbances) are along the same axis as the direction
of motion of the wave; transverse waves are at right angles to the direction of motion of
the wave. If this doesn't help, try imagining yourself as one of the particles that the wave
is disturbing (a water drop on the surface of the ocean, or an air molecule). As it comes
from behind you, a transverse waves lifts you up and then drops you down; a
longitudinal wave coming from behind pushes you forward and then pulls you back. You
can view animations of longitudinal and transverse waves here, single particles being
disturbed by a transverse wave or by a longitudinal wave, and particles being disturbed
by transverse and longitudinal waves. (There were also some nice animations of
longitudinal waves available as of this writing at Musemath.)

The Laplace Correction
An azimuth measured with a theodolite is an astronomic azimuth. It refers to an

azimuth expressed in a local horizon coordinate system aligned along the local gravity
vector (line of the plumb bob). A geodetic azimuth, on the other hand, is expressed
about a local normal to the ellipsoid. An astronomic azimuth is converted into a
geodetic azimuth by applying the Laplace correction. The Laplace correction typically
takes one of two forms:
The Simplified Laplace Correction (Hor. Laplace) :

- Eta x Tangent (Geodetic Latitude)
The Complete (Extended) Laplace Correction:

- Eta x Tangent(Geodetic Latitude)
- [Xi x Sine(Azimuth) - Eta x Cosine(Azimuth)] x Cotangent(Zenith Dist.)
The second term is known as the deflection correction, and it is negligible for horizontal
lines of sight.The output from DEFLEC provides not only the Xi and Eta components of
deflection of the vertical, but also provides the horizontal (simplified)Laplace correction.
The sense of the sign is such that one should ADD the correction from program
DEFLEC (Hor. Laplace) to an astronomic azimuth that is measured clockwise, thereby
obtaining a "Laplace Azimuth". Such an azimuth does not contain the deflection
correction, and is suitable for encoding into a *60* "Horizontal Blue-Book" record.
Effect Of Temperature
As you increase the temperature the rate of reaction increases. As a rough
approximation, for many reactions happening at around room temperature, the rate of
reaction doubles for every 10°C rise in temperature.
You have to be careful not to take this too literally. It doesn't apply to all reactions. Even
where it is approximately true, it may be that the rate doubles every 9°C or 11°C or
whatever. The number of degrees needed to double the ra Sound Propagation
Underwater

How does sound velocity vary in the oceans?

Since sound speed varies with temperature, pressure and salinity there are considerable
variations in sound velocity both spatially (with depth / geographically) and temporally
(daily / seasonally).
Horizontal variations in sound velocity are usually small due to small gradients in T, S
and p. Exceptions might occur in estuaries or around oceanic frontal systems. However,
since vertical gradients in T, p and S are much larger vertical variations in sound velocity
are much larger!
The simplified schematic above shows idealized vertical profiles of temperature and
sound velocity. This simplified diagram can be divided into three distinct zones. Closest
to the surface (zone 1) there is an isothermal layer created and maintained by mixing due
to wind and waves. Within this layer, which can be up to 200m deep the sound velocity,
increases slowly with depth due to the increasing pressure.

The middle layer (zone 2) is the thermocline. Here the sound velocity decreases rapidly
with depth due to the decreasing temperature. The base of the permanent thermocline
varies greatly with latitude but is typically found at a depth of about a 1000m.
Within the deepest region (zone 3) below the permanent thermocline the temperature
change is less dramatic. Here the sound velocity shows a further increase with depth due
to increasing pressure (like the surface layer).
Pressure and Salinity on Velocity Of Sound
Figure showing the convergence of sound rays around velocity minima (sound
channels).
Close examination of the figure above shows that the thermal structure within the water
column not only leads to a focusing acoustic energy but also can give rise to other areas
of acoustic darkness or 'shadow zones'. This also has important implications in
underwater acoustics. For example it would be impossible to detect an enemy submarine
with sonar if it were positioned in the shadow zone! Before the concept of shadow zones
was fully understood (before World War II) there was no explanation of the mysteriously
weak echoes that that were obtained in the afternoon when previously detection had been
good in the morning. At first sonar operators were blamed for their poor performance in
the afternoon. We now know that the weak echoes in the afternoon were due to the
development of the diurnal thermocline resulting from a warming of the water column
by the sun throughout the day. This phenomenon is now known as 'the afternoon effect'.

Sound channels may also exist in areas where there is an absence of any thermal
structure in the water column. For example in Polar regions where deep water masses are
formed (e.g. Antarctic Bottom Water and North Atlantic Deep Water) both salinity and
temperature change little with depth (isohaline and isothermal conditions). In these
regions sound velocity increases with depth due to the increasing pressure forming a
velocity minima at the surface. Thus, sound rays are continually refracted towards the
sea surface where they are strongly reflected. The process is continually repeated in a
similar way to the shallow water sound channel discussed earlier but on a larger scale.
This process is intensified in the presence of sea ice that leads to stronger reflections
from the sea surface.
In shallow water sound energy can also be focused even in the absence of a thermocline
due to repeated reflection from the sea bed and sea surface. Here shallow water refers to
the inland waters of bays, harbours and coastal waters less than 200m deep. In this type
of sound channel the characteristics of the sea bed and sea surface significantly affect the
efficiency of sound propagation. We will be examining this further when look at wave
reflection. Sound channels have often been used for research applications. For example
in the Mid-Ocean Dynamics Experiment (MODE) SOFAR floats were used to map
ocean currents. These were floats were designed to become negatively buoyant at the
depth of the sound channel. The floats emitted a low frequency signal at about 300Hz
that allowed them to be tracked by listening stations placed at known geographical
locations. The position of the floats was monitored as they drifted freely with the ocean
currents using the time of arrival of the acoustic pulses at the different locations. It was
devises such as these that initially demonstrated the existence of meso-scale eddies in the
ocean.
Examples
Some reactions are virtually instantaneous - for example, a precipitation reaction
involving the coming together of ions in solution to make an insoluble solid, or the
reaction between hydrogen ions from an acid and hydroxide ions from an alkali in
solution. So heating one of these won't make any noticeable difference to the rate of the
reaction.

Almost any other reaction you care to name will happen faster if you heat it - either in
the lab, or in industry.
The explanation
Particles can only react when they collide. If you heat a substance, the particles move
faster and so collide more frequently. That will speed up the rate of reaction.
That seems a fairly straightforward explanation until you look at the numbers!
It turns out that the frequency of two-particle collisions in gases is proportional to the
square root of the kelvin temperature. If you increase the temperature from 293 K to 303
K (20°C to 30°C), you will increase the collision frequency by a factor of:
That's an increase of 1.7% for a 10° rise. The rate of reaction will probably have doubled
for that increase in temperature - in other words, an increase of about 100%. The effect
of increasing collision frequency on the rate of the reaction is very minor. The important
effect is quite different . . .
The key importance of activation energy

Collisions only result in a reaction if the particles collide with enough energy to get the
reaction started. This minimum energy required is called the activation energy for the
reaction. Mark the position of activation energy on a Maxwell-Boltzmann distribution to
get a diagram like this:

Only those particles represented by the area to the right of the activation energy will
react when they collide. The great majority doesn’t have enough energy, and will simply
bounce apart.
To speed up the reaction, you need to increase the number of the very energetic particles
- those with energies equal to or greater than the activation energy. Increasing the
temperature has exactly that effect - it changes the shape of the graph.
In the next diagram, the graph labelled T is at the original temperature. The graph
labelled T+t is at a higher temperature.
If you now mark the position of the activation energy, you can see that although the
curve hasn't moved very much overall, there has been such a large increase in the
number of the very energetic particles that many more now collide with enough energy
to react

Remember that the area under a curve gives a count of the number of particles. On the
last diagram, the area under the higher temperature curve to the right of the activation
energy looks to have at least doubled - therefore at least doubling the rate of the reaction
Summary
Increasing the temperature increases reaction rates because of the disproportionately
large increase in the number of high energy collisions. It is only these collisions
(possessing at least the activation energy for the reaction) which result in a reaction.
Characteristics of Sound:
Intensity and the Decibel Scale
Sound waves are introduced into a medium by the vibration of an object. For example, a
vibrating guitar string forces surrounding air molecules to be compressed and expanded,
creating a pressure disturbance consisting of an alternating pattern of compressions and
rarefactions. The disturbance then travels from particle to particle through the medium,
transporting energy as it moves. The energy, which is carried by the disturbance, was
originally imparted to the medium by the vibrating string. The amount of energy which
is transferred to the medium is dependent upon the amplitude of vibrations of the guitar
string. If the more energy is put into the plucking of the string (that is, more work is done
to displace the string a greater amount from its rest position), then the string vibrates
with wider amplitude. The greater amplitude of vibration of the guitar string thus imparts
more energy to the medium, causing air particles to be displaced a greater distance from
their rest position. Subsequently, the amplitude of vibration of the particles of the
medium is increased, corresponding to an increased amount of energy being carried by

the particles. This relationship between energy and amplitude was discussed in more
detail in a previous unit.
The amount of energy, which is transported past a given area of the medium per unit of
time, is known as the intensity of the sound wave. The greater the amplitude of
vibrations of the particles of the medium, the greater the rate at which energy is
transported through it, and the more intense that the sound wave is. Intensity is the
energy/time/area; and since the energy/time ratio is equivalent to the quantity power,
intensity is simply the power/area.
Typical units for expressing the intensity of a sound wave are Watts/meter2.
As a sound wave carries its energy through a two-dimensional or three-dimensional

medium, the intensity of the sound wave decreases with increasing distance from the
source. The decrease in intensity with increasing distance is explained by the fact that
the wave is spreading out over a circular (2 dimensions) or spherical (3 dimensions)
surface and thus the energy of the sound wave is being distributed over a greater surface
area. The diagram at the right shows that the sound wave in a 2-dimensional medium is
spreading out in space over a circular pattern. Since energy is conserved and the area
through which this energy is transported is increasing, the power (being a quantity which
is measured on a per area basis) must decrease. The mathematical relationship between
intensity and distance is sometimes referred to as an inverse square relationship. As the
intensity varies inversely with the square of the distance from the source. So if the
distance from the source is doubled (increased by a factor of 2), then the intensity is
quartered (decreased by a factor of 4). Similarly, if the distance from the source is
quadrupled, then the intensity is decreased by a factor of 16. Applied to the diagram at
the right, the intensity at point B is one-fourth the intensity as point A and the intensity at
point C is one-sixteenth the intensity at point A. Since the intensity-distance relationship
is an inverse relationship, an increase in one quantity corresponds to a decrease in the
other quantity. And since the intensity-distance relationship is an inverse square
relationship, whatever factor by which the distance is increased, the intensity is

decreased by a factor equal to the square of the "distance change factor." The sample
data in the table below illustrate the inverse square relationship between power and
distance.
Distance Intensity
1m 160 units
2m 40 units
3m 17.8 units
4m 10 units
Humans are equipped with very sensitive ears capable of detecting sound waves of
extremely low intensity. The faintest sound that the typical human ear can detect has an
intensity of 1*10-12 W/m2. This intensity corresponds to a pressure wave in which a
compression of the particles of the medium increases the air pressure in that
compressional region by a mere 0.3 billionth of an atmosphere. A sound with an
intensity of 1*10-12 W/m2 corresponds to a sound, which will displace particles of air by
a mere one-billionth of a centimeter. The human ear can detect such a sound. WOW!
This faintest sound, which the human ear can detect, is known as the threshold of
hearing. The most intense sound, which the ear can safely detect without suffering any
physical damage, is more than one billion times more intense than the threshold of
hearing.
Since the range of intensities which the human ear can detect is so large, the scale which
is frequently used by physicists to measure intensity is a scale based on multiples of 10.
This type of scale is sometimes referred to as a logarithmic scale. The scale for
measuring intensity is the decibel scale. The threshold of hearing is assigned a sound
level of 0 decibels (abbreviated 0 dB); this sound corresponds to an intensity of 1*10 -12
W/m2. A sound which is 10 times more intense ( 1*10 -11 W/m2) is assigned a sound level
of 10 dB. A sound which is 10*10 or 100 times more intense ( 1*10 -10 W/m2) is assigned
a sound level of 20 db. A sound which is 10*10*10 or 1000 times more intense ( 1*10 -9
W/m2) is assigned a sound level of 30 db. A sound which is 10*10*10*10 or 10000 times
more intense ( 1*10-8 W/m2) is assigned a sound level of 40 db. Observe that this scale is
based on powers or multiples of 10. If one sound is 10x times more intense than another

sound, then it has a sound level which is 10*x more decibels than the less intense sound.
The table below lists some common sounds with an estimate of their intensity and
decibel level.
Intensity # of Times
Source Intensity
Level Greater Than TOH
-12 2
Threshold of Hearing (TOH) 1*10 W/m 0 dB 100
Rustling Leaves 1*10-11 W/m2 10 dB 101
Whisper 1*10-10 W/m2 20 dB 102
Normal Conversation 1*10-6 W/m2 60 dB 106
Busy Street Traffic 1*10-5 W/m2 70 dB 107
Vacuum Cleaner 1*10-4 W/m2 80 dB 108
Large Orchestra 6.3*10-3 W/m2 98 dB 109.8
Walkman at Maximum Level 1*10-2 W/m2 100 dB 1010
Front Rows of Rock Concert 1*10-1 W/m2 110 dB 1011
Threshold of Pain 1*101 W/m2 130 dB 1013
Military Jet Takeoff 1*102 W/m2 140 dB 1014
Instant Perforation of Eardrum 1*104 W/m2 160 dB 1016
Loudness
While the intensity of a sound is a very objective quantity which can be measured with
sensitive instrumentation, the loudness of a sound is more of a subjective response
which will vary with a number of factors. The same sound will not be perceived to have
the same loudness to all individuals. Age is one factor which effects the human ear's
response to a sound. Quite obviously, your grandparents do not hear like they used to.
The same intensity sound would not be perceived to have the same loudness to them as it
would to you. Furthermore, two sounds with the same intensity but different frequencies
will not be perceived to have the same loudness. Because of the human ear's tendency to
amplify sounds having frequencies in the range from 1000 Hz to 5000 Hz, sounds with
these intensities seem louder to the human ear. Despite the distinction between intensity
and loudness, it is safe to state that the more intense sounds will be perceived to be the
loudest sounds
Pitch and Frequency

A sound wave, like any other wave, is introduced into a medium by a vibrating object.
The vibrating object is the source of the disturbance which moves through the medium.
The vibrating object which creates the disturbance could be the vocal chords of a person,
the vibrating string and sound board of a guitar or violin, the vibrating tines of a tuning
fork, or the vibrating diaphragm of a radio speaker. Regardless of what vibrating object
is creating the sound wave, the particles of the medium through which the sound moves
is vibrating in a back and forth motion at a given frequency. The frequency of a wave
refers to how often the particles of the medium vibrate when a wave passes through the
medium. The frequency of a wave is measured as the number of complete back-and-
forth vibrations of a particle of the medium per unit of time. If a particle of air undergoes
1000 longitudinal vibrations in 2 seconds, then the frequency of the wave would be 500
vibrations per second. A commonly used unit for frequency is the Hertz (abbreviate Hz)
where 1 Hertz = 1 vibration/second
As a sound wave moves through a medium, each particle of the medium vibrates at the
same frequency. This is sensible since each particle vibrates due to the motion of its
nearest neighbor. The first particle of the medium begins vibrating, at say 500 Hz, and
begins to set the second particle into vibrational motion at the same frequency of 500
Hz. The second particle begins vibrating at 500 Hz and thus sets the third particle of the
medium into vibrational motion at 500 Hz. The process continues throughout the
medium; each particle vibrates at the same frequency. And of course the frequency at
which each particle vibrates is the same as the frequency of the original source of the
sound wave. Subsequently, a guitar string vibrating at 500 Hz will set the air particles in
the room vibrating at the same frequency of 500 Hz which carries a sound signal to the
ear of a listener which is detected as a 500 Hz sound wave.
The back-and-forth vibration motion of the particles of the medium would not be the
only observable phenomenon occurring at a given frequency. Since a sound wave is a
pressure wave, a detector could be used to detect oscillations in pressure from a high
pressure to a low pressure and back to a high pressure. As the compression (high
pressure) and rarefaction (low pressure) disturbances move through the medium, they
would reach the detector at a given frequency. For example, a compression would reach
the detector 500 times per second if the frequency of the wave were 500 Hz. Similarly, a

rarefaction would reach the detector 500 times per second if the frequency of the wave
were 500 Hz. Thus the frequency of a sound wave not only refers to the number of back-
and-forth vibrations of the particles per unit of time, but also refers to the number of
compression or rarefaction disturbances which pass a given point per unit of time. A
detector could be used to detect the frequency of these pressure oscillations over a given
period of time. The typical output provided by such a detector is a pressure-time plot as
shown below.
Since a pressure-time plot shows the fluctuations in pressure over time, the period of the
sound wave can be found by measuring the time between successive high pressure points
(corresponding to the compressions) or the time between successive low pressure points
(corresponding to the rarefactions). As discussed in an earlier unit, the frequency is
simply the reciprocal of the period. For this reason, a sound wave with a high frequency
would correspond to a pressure time plot with a small period - that is, a plot
corresponding to a small amount of time between successive high pressure points.
Conversely, a sound wave with a low frequency would correspond to a pressure time
plot with a large period - that is, a plot corresponding to a large amount of time between
successive high pressure points. The diagram below shows two pressure-time plots, one
corresponding to a high frequency and the other to a low frequency.

The ears of humans (and other animals) are sensitive detectors capable of detecting the
fluctuations in air pressure, which impinge upon the eardrum. The mechanics of the ear's
detection ability will be discussed later in this lesson. For now, it is sufficient to say that
the human ear is capable of detecting sound waves with a wide range of frequencies,
ranging between approximately 20 Hz to 20 000 Hz. Any sound with a frequency below
the audible range of hearing (i.e., less than 20 Hz) is known as an infrasound and any
sound with a frequency above the audible range of hearing (i.e., more than 20 000 Hz) is
known as an ultrasound. Humans are not alone in their ability to detect a wide range of
frequencies. Dogs can detect frequencies as low as approximately 50 Hz and as high as
45 000 Hz. Cats can detect frequencies as low as approximately 45 Hz and as high as 85
000 Hz. Bats, who are essentially blind and must rely on sound echolocation for
navigation and hunting, can detect frequencies as high as 120 000 Hz. Dolphins can
detect frequencies as high as 200 000 Hz. While dogs, cats, bats, and dolphins have an
unusual ability to detect ultrasound, an elephant possesses the unusual ability to detect
infrasound, having an audible range from approximately 5 Hz to approximately 10 000
Hz.
The sensations of these frequencies are commonly referred to as the pitch of a
sound. A high pitch sound corresponds to a high frequency and a low pitch sound
corresponds to a low frequency. Amazingly, many people, especially those who have
been musically trained, are capable of detecting a difference in frequency between
two separate sounds which is as little as 2 Hz. When two sounds with a frequency
difference of greater than 7 Hz are played simultaneously, most people are capable of

detecting the presence of a complex wave pattern resulting from the interference and
superposition of the two sound waves. Certain sound waves when played (and heard)
simultaneously will produce a particularly pleasant sensation when heard, are are
said to be consonant. Such sound waves form the basis of intervals in music. For
example, any two sounds whose frequencies make a 2:1 ratio are said to be separated
by an octave and result in a particularly pleasing sensation when heard; that is, two
sound waves sound good when played together if one sound has twice the frequency
of the other. Similarly two sounds with a frequency ratio of 5:4 are said to be
separated by an interval of a third; such sound waves also sound good when played
together. Examples of other sound wave intervals and their respective frequency
ratios are listed in the table below.
Interval Frequency Ratio Examples

Octave 2:1 512 Hz and 256 Hz
Third 5:4 320 Hz and 256 Hz
Fourth 4:3 342 Hz and 256 Hz
Fifth 3:2 384 Hz and 256 Hz
The ability of humans to perceive pitch is associated with the frequency of the
sound wave, which impinges upon the ear. Because sound waves are longitudinal
waves, which produce high- and low-pressure disturbances of the particles of a
medium at a given frequency, the ear has an ability to detect such frequencies and
associate them with the pitch of the sound. But pitch is not the only property of a
sound wave detectable by the human ear
1. Two notes, which have a frequency ratio of 2:1, are said to be separated by an
octave. A note which is separated by an octave from middle C (256 Hz) is
a) 128 Hz b) 254 Hz c) 258 Hz d) 345 Hz e) none of these
SIREN:

It consists of a wind chest. The upper disc has holes near its circumference and is just
above the disc A of the wind chest. The holes in D are slanting and are just in the
opposite direction to the holes in the fixed disc A. the disc D rotates about the spindle
and the number of rotations made by it, is indicated by the counters C1 and C2.
Working: the wind is pushed into the wind chest and each puff of air that escapes from
the holes, rotates the upper disc D. when the holes in the Disc D are just above the holes
in disc A, air escapes. Suppose the upper disc has h holes and it rotates n times second.
The frequency of the sound produced by the air escaping out of the holes =N= n x h
Determination of the frequency of a note:

The pressure of wind entering into the wind chest is added is adjusted such that
frequency of the sound produced by the siren is nearly equal to that given note.
Increasing the speed of the rotation of the disc D slowly can do this. When the
frequencies are nearly equal, count the number of beats produced per second. Let the
number of beats produced per second be x.
The frequency of the note = N + x
To avoid the counting of the beats, the speed of the rotation should be adjusted such that
the two notes have the same pitch. In that case the frequency of the note = N
Example: the disc of a given siren has 32 holes. What must be the speed of the rotation
per minute of the siren disc so that the note emitted by the siren may be in unison with
that of the tuning fork, which makes 256 vibrations per second?
Suppose, the number of rotations/minute = x
Therefore Number of rotations = x/60
Here n = x/60
h= 32, N = 256
N = n* h
256 = x/60 *32
x = 480 rotations/minute.

Doppler’s Effect:
It is commonly observed that the pitch of a note apparently changes when either the
source or the observer are in motion relative to each other. When the source approaches
the observer or when the observer approaches the source or when both approach each
other the apparent pitch is higher than the actual pitch of the sound produced by the
source.
Suppose a person is standing on the platform. the apparent pitch of the whistle of the
engine increases, when the engine is approaching the person. When the engine moves
away from the person, the apparent pitch of the whistle of the engine decreases. This
apparent change in the pitch due to relative motion between the source and the observer
is called the Doppler’s Effect.
The Doppler’s Effect occurs when a source or light moves relative to an observer. In
light, the effect was observed when measurements were taken of the wavelength of the
light from a moving star; they showed a mark variation. In sound, this effect can be
demonstrated by placing a whistle in the end of a long piece of rubber tubing, and
whirling the tube in a horizontal circle above the head while blowing the whistle. The
open end of the rubber tube acts as a moving source of sound, and an observer hears a
rise and fall in the pitch as the end approaches and recedes from him or her.
Example: Police car sirens use high frequency sound. When a police car is chasing a
speeding car, the sound moving towards you appears to increase in frequency or pitch.
When the car passes and moves away, the frequency appears to be lower. As we see
soon, the frequency change is due to the Dopplers effect. A similar change occurs when
Ambulance sirens are sounding and the ambulance is moving at high speed towards or
away from you.

Applications
1) Velocity and rotation of the sun:
By the study of the Doppler’s shift from the light received from the western and
the eastern edges of the sun, it has been found that the shift is due to the velocity of
about 2Kms. Moreover, no such shift has been observed from the light received from
the north and south edges. This show that the sun rotates about the north south axis.
2) Discovery of double stars:
By the constant observation of the sky, it has been foud that some of the stars that
appear to be single are actually double stars and is known as spectroscopic binaries.
These stars revolve about each other. When one is Approaching the earth, and the other
is going away from the earth, there is a shift in their spectral lines and a single spectral
line is split up into two lines whose separation depends upon time and the time period
is equal to the time period of revolution of the stars. By this method a number of
double star systems have been found.
3) Red shift:
It has been observed that some distant nebulae are moving away with the velocity
greater than 20 x 10 raise to 3 Kms.and the important spectral lines appear to shift
towards the red end of the spectrum by 200 A. this gives the idea that the universe is
expanding.
4) Saturn’s ring:
The planet Saturn has been found to be surrounded by concentric rings. With the
help of Doppler’s effect it has been found that these rings are not solids but consists of
a number of ‘satellites’ moving around the Saturn in these orbits. If the rings were
solids, the outer edge of the ring should have greater velocity that the inner edge. But
according to the principle of a satellite the velocity of the satellite in the inner orbit is
more than the outer orbit. This fact has been established by the Dopplers shift. Thus,
the rings of the Saturn are not solids but they are large numbers of satellites moving in
these orbits called the rings of the Saturn.

6. Light
Propagation of light
The sun radiates light and heat energy. The light energy are transported in the form of
rays. People usually get the impression that we can 'see' rays. This phenomenon is
caused by the scattering of light by the suspended dust particles or water droplets.
Scientists actually invented the idea of rays. Thus, a ray of line is actually referring to the
line along which the light travels. On ray diagrams, they are actually represented by lines
with arrows. In reality, they have no arrows but they are used to indicate the direction of
the ight rays. A bundle of rays is known as a beam of light. The figure below shows you
the 3 types of light beams.
So much for theory, now let's do the following experiment for you to see on your own
that light really travels in straight lines
The apparent incompatibility of the law of Propagation of light with the principle
of relativity
There is hardly a simpler law in physics than that according to which light is propagated
in empty space. Every child at school knows, or believes he knows, that this propagation
takes place in straight lines with a velocity c= 300,000 km./sec. At all events we know
with great exactness that this velocity is the same for all colours, because if this were not
the case, the minimum of emission would not be observed simultaneously for different
colours during the eclipse of a fixed star by its dark neighbour. By means of similar
considerations based on observations of double stars, the Dutch astronomer De Sitter
was also able to show that the velocity of propagation of light cannot depend on the
velocity of motion of the body emitting the light. The assumption that this velocity of
propagation is dependent on the direction "in space" is in itself improbable.
In short, let us assume that the simple law of the constancy of the velocity of light c (in
vacuum) is justifiably believed by the child at school. Who would imagine that this
simple law has plunged the conscientiously thoughtful physicist into the greatest
intellectual difficulties? Let us consider how these difficulties arise.
Of course we must refer the process of the propagation of light (and indeed every other
process) to a rigid reference-body (co-ordinate system). As such a system let us again

choose our embankment. We shall imagine the air above it to have been removed. If a
ray of light be sent along the embankment, we see from the above that the tip of the ray
will be transmitted with the velocity c relative to the embankment. Now let us suppose
that our railway carriage is again travelling along the railway lines with the velocity v,
and that its direction is the same as that of the ray of light, but its velocity of course
much less. Let us inquire about the velocity of propagation of the ray of light relative to
the carriage. It is obvious that we can here apply the consideration of the previous
section, since the ray of light plays the part of the man walking along relatively to the
carriage. The velocity w of the man relative to the embankment is here replaced by the
velocity of light relative to the embankment. w is the required velocity of light with
respect to the carriage, and we have
w = c-v.
The velocity of propagation ot a ray of light relative to the carriage thus comes cut
smaller than c.
But this result comes into conflict with the principle of relativity set forth in Section V.
For, like every other general law of nature, the law of the transmission of light in vacuo
[in vacuum] must, according to the principle of relativity, be the same for the railway
carriage as reference-body as when the rails are the body of reference. But, from our
above consideration, this would appear to be impossible. If every ray of light is
propagated relative to the embankment with the velocity c, then for this reason it would
appear that another law of propagation of light must necessarily hold with respect to the
carriage — a result contradictory to the principle of relativity.
In view of this dilemma there appears to be nothing else for it than to abandon either the
principle of relativity or the simple law of the propagation of light in vacuo. Those of
you who have carefully followed the preceding discussion are almost sure to expect that
we should retain the principle of relativity, which appeals so convincingly to the intellect
because it is so natural and simple. The law of the propagation of light in vacuo would
then have to be replaced by a more complicated law conformable to the principle of
relativity. The development of theoretical physics shows, however, that we cannot pursue
this course. The epoch-making theoretical investigations of H. A. Lorentz on the
electrodynamical and optical phenomena connected with moving bodies show that

experience in this domain leads conclusively to a theory of electromagnetic phenomena,

of which the law of the constancy of the velocity of light in vacuo is a necessary
consequence. Prominent theoretical physicists were theref ore more inclined to reject the
principle of relativity, in spite of the fact that no empirical data had been found which
were contradictory to this principle.
At this juncture the theory of relativity entered the arena. As a result of an analysis of the
physical conceptions of time and space, it became evident that in realily there is not the
least incompatibilitiy between the principle of relativity and the law of propagation of
light, and that by systematically holding fast to both these laws a logically rigid theory
could be arrived at. This theory has been called the special theory of relativity to
distinguish it from the extended theory, with which we shall deal later. In the following
pages we shall present the fundamental ideas of the special theory of relativity
Visible light is a narrow part of the electromagnetic spectrum and in a vacuum all
electromagnetic radiation travels at the speed of light:
The above number is now accepted as a standard value and the value of the meter is
defined to be consistent with it. In a material medium the effective speed of light is
slower and is usually stated in terms of the index of refraction of the medium. Light
propagation is affected by the phenomena refraction, reflection, diffraction, and
interference.
The behavior of light in optical systems will be characterized in terms of its vergence.

The reflection and refraction of light
Reflections
A reflection flips all the points in the plane over a line, which is called the mirror. A
reflection changes the the sense of figures in the plane. All the points in the mirror
contains all the invariant points (points that map onto themselves) under a reflection.
Rays and wave fronts
Light is a very complex phenomenon, but in many situations its behavior can be
understood with a simple model based on rays and wave fronts. A ray is a thin beam of
light that travels in a straight line. A wave front is the line (not necessarily straight) or
surface connecting all the light that left a source at the same time. For a source like the
Sun, rays radiate out in all directions; the wave fronts are spheres centered on the Sun. If
the source is a long way away, the wave fronts can be treated as parallel lines.
Rays and wave fronts can generally be used to represent light when the light is
interacting with objects that are much larger than the wavelength of light, which is about
500 nm. In particular, we'll use rays and wave fronts to analyze how light interacts with
mirrors and lenses.
The law of reflection
Objects can be seen by the light they emit, or, more often, by the light they reflect.
Reflected light obeys the law of reflection, that the angle of reflection equals the angle of
incidence.
For objects such as mirrors, with surfaces so smooth that any hills or valleys on the
surface are smaller than the wavelength of light, the law of reflection applies on a large
scale. All the light travelling in one direction and reflecting from the mirror is reflected
in one direction; reflection from such objects is known as specular reflection.

Most objects exhibit diffuse reflection, with light being reflected in all directions. All
objects obey the law of reflection on a microscopic level, but if the irregularities on the
surface of an object are larger than the wavelength of light, which is usually the case, the
light reflects off in all directions.
Plane mirrors
A plane mirror is simply a mirror with a flat surface; all of us use plane mirrors every
day, so we've got plenty of experience with them. Images produced by plane mirrors
have a number of properties, including:
the image produced is upright
the image is the same size as the object (i.e., the magnification is m = 1)
the image is the same distance from the mirror as the object appears to be (i.e., the image
distance = the object distance)
the image is a virtual image, as opposed to a real image, because the light rays do not
actually pass through the image. This also implies that an image could not be focused on
a screen placed at the location where the image is.
A little geometry
Dealing with light in terms of rays is known as geometrical optics, for good reason: there
is a lot of geometry involved. It's relatively straight-forward geometry, all based on
similar triangles, but we should review that for a plane mirror.
Consider an object placed a certain distance in front of a mirror, as shown in the
diagram. To figure out where the image of this object is located, a ray diagram can be
used. In a ray diagram, rays of light are drawn from the object to the mirror, along with
the rays that reflect off the mirror. The image will be found where the reflected rays
intersect. Note that the reflected rays obey the law of reflection. What you notice is that
the reflected rays diverge from the mirror; they must be extended back to find the place
where they intersect, and that's where the image is.

Analyzing this a little further, it's easy to see that the height of the image is the same as
the height of the object. Using the similar triangles ABC and EDC, it can also be seen
that the distance from the object to the mirror is the same as the distance from the image
to the mirror.
Spherical mirrors
Light reflecting off a flat mirror is one thing, but what happens when light reflects off a
curved surface? We'll take a look at what happens when light reflects from a spherical
mirror, because it turns out that, using reasonable approximations, this analysis is fairly
straight-forward. The image you see is located either where the reflected light converges,
or where the reflected light appears to diverge from.
A spherical mirror is simply a piece cut out of a reflective sphere. It has a center of
curvature, C, which corresponds to the center of the sphere it was cut from; a radius of
curvature, R, which corresponds to the radius of the sphere; and a focal point (the point
where parallel light rays are focused to) which is located half the distance from the
mirror to the center of curvature. The focal length, f, is therefore:
focal length of a spherical mirror : f = R / 2
This is actually an approximation. Parabolic mirrors are really the only mirrors that
focus parallel rays to a single focal point, but as long as the rays don't get too far from
the principal axis then the equation above applies for spherical mirrors. The diagram
shows the principal axis, focal point (F), and center of curvature for both a concave and
convex spherical mirror.

Spherical mirrors are either concave (converging) mirrors or convex (diverging) mirrors,
depending on which side of the spherical surface is reflective. If the inside surface is
reflective, the mirror is concave; if the outside is reflective, it's a convex mirror. Concave
mirrors can form either real or virtual images, depending on where the object is relative
to the focal point. A convex mirror can only form virtual images. A real image is an
image that the light rays from the object actually pass through; a virtual image is formed
because the light rays can be extended back to meet at the image position, but they don't
actually go through the image position.
Ray diagrams
To determine where the image is, it is very helpful to draw a ray diagram. The image will
be located at the place where the rays intersect. You could just draw random rays from
the object to the mirror and follow the reflected rays, but there are three rays in particular
that are very easy to draw.
Only two rays are necessary to locate the image on a ray diagram, but it's useful to add
the third as a check. The first is the parallel ray; it is drawn from the tip of the object
parallel to the principal axis. It then reflects off the mirror and either passes through the
focal point, or can be extended back to pass through the focal point.
The second ray is the chief ray. This is drawn from the tip of the object to the mirror
through the center of curvature. This ray will hit the mirror at a 90° angle, reflecting
back the way it came. The chief and parallel rays meet at the tip of the image.
The third ray, the focal ray, is a mirror image of the parallel ray. The focal ray is drawn
from the tip of the object through (or towards) the focal point, reflecting off the mirror
parallel to the principal axis. All three rays should meet at the same point.

A ray diagram for a concave mirror is shown above. This shows a few different things.
For this object, located beyond the center of curvature from the mirror, the image lies
between the focal point (F) and the center of curvature. The image is inverted compared
to the object, and it is also a real image, because the light rays actually pass through the
point where the image is located.
With a concave mirror, any object beyond C will always have an image that is real,
inverted compared to the object, and between F and C. You can always trade the object
and image places (that just reverses all the arrows on the ray diagram), so any object
placed between F and C will have an image that is real, inverted, and beyond C. What
happens when the object is between F and the mirror? You should draw the ray diagram
to convince yourself that the image will be behind the mirror, making it a virtual image,
and it will be upright compared to the object.
A ray diagram for a convex mirror
What happens with a convex mirror? In this case the ray diagram looks like this:
As the ray diagram shows, the image for a convex mirror is virtual, and upright
compared to the object. A convex mirror is the kind of mirror used for security in stores,
and is also the kind of mirror used on the passenger side of many cars ("Objects in
mirror are closer than they appear."). A convex mirror will reflect a set of parallel rays in

all directions; conversely, it will also take light from all directions and reflect it in one
direction, which is exactly how it's used in stores and cars.
The mirror equation
Drawing a ray diagram is a great way to get a rough idea of how big the image of an
object is, and where the image is located. We can also calculate these things precisely,
using something known as the mirror equation. The textbook does a nice job of deriving
this equation in section 25.6, using the geometry of similar triangles.
Magnification
In most cases the height of the image differs from the height of the object, meaning that
the mirror has done some magnifying (or reducing). The magnification, m, is defined as
the ratio of the image height to the object height, which is closely related to the ratio of
the image distance to the object distance:
A magnification of 1 (plus or minus) means that the image is the same size as the object.
If m has a magnitude greater than 1 the image is larger than the object, and an m with a
magnitude less than 1 means the image is smaller than the object. If the magnification is
positive, the image is upright compared to the object; if m is negative; the image is
inverted compared to the object.
Sign conventions
What does a positive or negative image height or image distance mean? To figure out
what the signs mean, take the side of the mirror where the object is to be the positive
side. Any distances measured on that side are positive. Distances measured on the other
side are negative.
f, the focal length, is positive for a concave mirror, and negative for a convex mirror.
When the image distance is positive, the image is on the same side of the mirror as the
object, and it is real and inverted. When the image distance is negative, the image is
behind the mirror, so the image is virtual and upright.
A negative m means that the image is inverted. Positive means an upright image.
Steps for analyzing mirror problems

There are basically three steps to follow to analyze any mirror problem, which generally
means determining where the image of an object is located, and determining what kind
of image it is (real or virtual, upright or inverted).
Step 1 - Draw a ray diagram. The more careful you are in constructing this, the better
idea you'll have of where the image is.
Step 2 - Apply the mirror equation to determine the image distance. (Or to find the
object distance, or the focal length, depending on what is given.)
Step 3 - Make sure steps 1 and 2 are consistent with each other.
An example
A Star Wars action figure, 8.0 cm tall, is placed 23.0 cm in front of a concave mirror
with a focal length of 10.0 cm. Where is the image? How tall is the image? What are the
characteristics of the image?
The first step is to draw the ray diagram, which should tell you that the image is real,
inverted, smaller than the object, and between the focal point and the center of curvature.
The location of the image can be found from the mirror equation:
which can be rearranged to:
The image distance is positive, meaning that it is on the same side of the mirror as the
object. This agrees with the ray diagram. Note that we don't need to worry about
converting distances to meters; just make sure everything has the same units, and
whatever unit goes into the equation is what comes out.
Calculating the magnification gives:
Solving for the image height gives:
The negative sign for the magnification, and the image height, tells us that the image is
inverted compared to the object.
To summarize, the image is real, inverted, 6.2 cm high, and 17.7 cm in front of the
mirror.

Example 2 - a convex mirror

The same Star Wars action figure, 8.0 cm tall, is placed 6.0 cm in front of a convex
mirror with a focal length of -12.0 cm. Where is the image in this case, and what are the
image characteristics?
Again, the first step is to draw a ray diagram. This should tell you that the image is
located behind the mirror; that it is an upright, virtual image; that it is a little smaller
than the object; and that the image is between the mirror and the focal point.
The second step is to confirm all those observations. The mirror equation, rearranged as
in the first example, gives:
Solving for the magnification gives:
This gives an image height of 0.667 x 8 = 5.3 cm.

All of these results are consistent with the conclusions drawn from the ray diagram. The
image is 5.3 cm high, virtual, upright compared to the object, and 4.0 cm behind the
mirror.
Refraction
When we talk about the speed of light, we're usually talking about the speed of light in a
vacuum, which is 3.00 x 108 m/s. When light travels through something else, such as
glass, diamond, or plastic, it travels at a different speed. The speed of light in a given
material is related to a quantity called the index of refraction, n, which is defined as the
ratio of the speed of light in vacuum to the speed of light in the medium:
index of refraction : n = c / v
When light travels from one medium to another, the speed changes, as does the
wavelength. The index of refraction can also be stated in terms of wavelength:
Although the speed changes and wavelength changes, the frequency of the light will be
constant. The frequency, wavelength, and speed are related by:

The change in speed that occurs when light passes from one medium to another is
responsible for the bending of light, or refraction, that takes place at an interface. If light
is travelling from medium 1 into medium 2, and angles are measured from the normal to
the interface, the angle of transmission of the light into the second medium is related to
the angle of incidence by Snell's law :
Solar Eclipses
One consequence of the Moon's orbit about the Earth is that the Moon can shadow the
Sun's light as viewed from the Earth, or the Moon can pass through the shadow cast by
the Earth. The former is called a solar eclipse and the later is called a lunar eclipse. The
small tilt of the Moon's orbit with respect to the plane of the ecliptic and the small
eccentricity of the lunar orbit make such eclipses much less common than they would be
otherwise, but partial or total eclipses are actually rather frequent.
Frequency of Eclipses
For example there will be 18 solar eclipses from 1996-2020 for which the eclipse will be
total on some part of the Earth's surface. The common perception that eclipses are
infrequent is because the observation of a total eclipse from a given point on the
surface of the Earth is not a common occurrence. For example, it will be two decades
before the next total solar eclipse visible in North America occurs.
The next total solar eclipse will be on August 11, 1999, with the path of totality crossing
the North Atlantic, Europe, the Middle East, and India. In this section we consider solar
eclipses and in the next we discuss lunar eclipses.
Geometry of Solar Eclipses

The geometry associated with solar eclipses is illustrated in the following figure (which,
like most figures in this and the next section, is illustrative and not to scale).
Geometry of solar eclipses
The shadow cast by the Moon can be divided by geometry into the completely shadowed
umbra and the partially shadowed penumbra.
Types of Solar Eclipses
The preceding figure allows three general classes of solar eclipses (as observed from any
particular point on the Earth) to be defined:
Total Solar Eclipses occur when the umbra of the Moon's shadow touches a region on
the surface of the Earth.
Partial Solar Eclipses occur when the penumbra of the Moon's shadow passes over a
region on the Earth's surface.
Annular Solar Eclipses occur when a region on the Earth's surface is in line with the
umbra, but the distances are such that the tip of the umbra does not reach the Earth's
surface.
As illustrated in the figure, in a total eclipse the surface of the Sun is completely blocked
by the Moon, in a partial eclipse it is only partially blocked, and in an annular eclipse the
eclipse is partial, but such that the apparent diameter of the Moon can be seen
completely against the (larger) apparent diameter of the Sun.
A given solar eclipse may be all three of the above for different observers. For example,
in the path of totality (the track of the umbra on the Earth's surface) the eclipse will be
total, in a band on either side of the path of totality the shadow cast by the penumbra

leads to a partial eclipse, and in some eclipses the path of totality extends into a path
associated with an annular eclipse because for that part of the path the umbra does not
reach the Earth's surface.
Total Solar Eclipses
A total solar eclipse requires the umbra of the Moon's shadow to touch the surface of the
Earth. Because of the relative sizes of the Moon and Sun and their relative distances
from Earth, the path of totality is usually very narrow (hundreds of kilometers across).
The following figure illustrates the path of totality produced by the umbra of the Moon's
shadow. (We do not show the penumbra, which will produce a partial eclipse in a much
larger region on either side of the path of totality; we also illustrate in this figure the
umbra of the Earth's shadow, which will be responsible for total lunar eclipses to be
discussed in the next section.)
Solar eclipse (not to scale)
As noted above, the images that we show in discussing eclipses are illustrative but not
drawn to scale. The true relative sizes of the Sun and Earth and Moon, and their
distances, are very different than in the above figure.
Animations of Solar Eclipses

Here are three animations that illustrate observations in a solar eclipse. The first
demonstrates generally the case of a total solar eclipse; the next two are simulated views

of two recent solar eclipses from unusual vantage points, one from the Moon and one
from the Sun (these last two were constructed using the program Starry Night).
Appearance of a Total Solar Eclipse

If you are in the path of totality the eclipse begins with a partial phase in which the
Moon gradually covers more and more of the Sun. This typically lasts for about an hour
until the Moon completely covers the Sun and the total eclipse begins. The duration of
totality can be as short as a few seconds, or as long as about 8 minutes, depending on the
details.
As totality approaches the sky becomes dark and a twilight that can only be described as
eerie begins to descend. Just before totality waves of shadow rushing rapidly from
horizon to horizon may be visible. In the final instants before totality light shining
through valleys in the Moon's surface gives the impression of beads on the periphery of
the Moon (a phenomenon called Bailey's Beads). The last flash of light from the surface
of the Sun as it disappears from view behind the Moon gives the appearance of a
diamond ring and is called, appropriately, the diamond ring effect (image at right).
As totality begins , the solar corona (extended outer atmosphere of the Sun) blazes into
view. The corona is a million times fainter than the surface of the Sun; thus only when
the eclipse is total can it be seen; if even a tiny fraction of the solar surface is still visible
it drowns out the light of the corona. At this point the sky is sufficiently dark that planets
and brighter stars are visible, and if the Sun is active one can typically see solar
prominences and flares around the limb of the Moon, even without a telescope (see
image at left).
The period of totality ends when the motion of the Moon begins to uncover the surface
of the Sun, and the eclipse proceeds through partial phases for approximately an hour
until the Sun is once again completely uncovered. Here is a movie of the 1994 total solar
eclipse (3.1 MB MPEG; Source; here is a QuickTime version, but note that it is 15 MB
in length).

A partial solar eclipse is interesting; a total solar eclipse is awe-inspiring in the literal
meaning of the phrase. If you have an opportunity to observe a total solar eclipse, don't
miss it! It is an experience that you will never forget.
Patterns of Eclipses
Because solar eclipses are the result of periodic motion of the Moon about the Earth,
there are regularities in the timing of eclipses that give cycles of related eclipses. These
cycles were known and used to predict eclipses long before there was a detailed
scientific understanding of what causes eclipses. For example, the ancient Babylonians
understood one such set of cycles called the Saros, and were able to predict eclipses
based on this knowledge. Here is a link to a discussion of such cycles and regularities in
eclipse patterns
Phases of the total lunar eclipse 1993 Nov 29.
What is an eclipse of the Moon? What causes eclipses and why? How often do eclipses
happen and when is the next eclipse of the Moon? You'll learn the answers to these
questions and more in MrEclipse's primer on lunar eclipses.
The Moon is a cold, rocky body about 2,160 miles (3,476 km) in diameter. It has no light
of its own but shines by sunlight reflected from its surface. The Moon orbits Earth about
once every 29 and a half days. As it circles our planet, the changing position of the Moon
with respect to the Sun causes our natural satellite to cycle through a series of phases:
New, New Crescent, First Quarter, Waxing Gibbous, Full,
Waning Gibbous, Last Quarter, Old Crescent and back to New again.
Phases of theMoon.

The phase known as New Moon can not actually be seen because the illuminated side of
the Moon is then pointed away from Earth. The rest of the phases are familiar to all of us
as the Moon cycles through them month after month. Did you realize that the word
month is derived from the Moon's 29.5 day period?
To many of us, Full Moon is the phase of love and romance. When the Moon is Full, it
rises at sunset and is visible all night long. At the end of the night, the Full Moon sets
just as the Sun rises. None of the Moon's other phases have this unique characteristic. It
happens because the Moon is directly opposite the Sun in the sky when the Moon is Full.
Full Moon also has special significance with regard to eclipses.
Geometry of the Sun, Earth and Moon During an Eclipse of the Moon.
Earth's two shadows are the penumbra and the umbra.
(Sizes and distances not to scale)
An eclipse of the Moon (or lunar eclipse) can only occur at Full Moon, and only if the
Moon passes through some portion of the Earth's shadow. The shadow is actually
composed of two cone-shaped components, one nested inside the other. The outer or
penumbral shadow is a zone where the Earth blocks part but not all of the Sun's rays

from reaching the Moon. In contrast, the inner or umbral shadow is a region where the
Earth blocks all direct sunlight from reaching the Moon.
Astronomers recognize three basic types of lunar eclipses:
1. Penumbral Lunar Eclipse
The Moon passes through Earth's penumbral shadow.
These events are of only academic interest since they are subtle and quite difficult to
observe.
2. Partial Lunar Eclipse
A portion of the Moon passes through Earth's umbral shadow.
These events are easy to see, even with the unaided eye.
3. Total Lunar Eclipse
The entire Moon passes through Earth's umbral shadow.
These events are quite striking for the vibrant range of colors the Moon can take on
during the total phase (i.e. - totality).
Now you might be wondering "If the Moon orbits Earth every 29.5 days and lunar
eclipses only occur at Full Moon, then why don't we have an eclipse once a month
during Full Moon?". I'm glad you asked! You see, the Moon's orbit around Earth is
actually tipped about 5 degrees to Earth's orbit around the Sun. This means that the
Moon spends most of the time either above or below the plane of Earth's orbit. And the
plane of Earth's orbit around the Sun is important because Earth's shadows lie exactly in
the same plane. During Full Moon, our natural satellite usually passes above or below
Earth's shadows and misses them entirely. No eclipse takes place. But two to four times
each year, the Moon passes through some portion of the Earth's penumbral or umbral
shadows and one of the above three types of eclipses occurs.
When an eclipse of the Moon takes place, everyone on the night side of Earth can see it.
About 35% of all eclipses are of the penumbral type which are very difficult to detect,
even with a telescope. Another 30% are partial eclipses which are easy to see with the
unaided eye. The final 35% or so are total eclipses, and these are quite extrordinary
events to behold.

Beginning (right), middle (center) and end (left) of totality

during the total lunar eclipse of 2000 Jan 20-21.
During a total lunar eclipse, the Earth blocks the Sun's light from reaching the Moon.
Astronauts on the Moon would then see the Earth eclipsing the Sun. (They would see a
bright red ring around the Earth as they watched all the sunrises and sunsets happening
simultaneousely around the world!) While the Moon remains completely within Earth's
umbral shadow, indirect sunlight still manages to reach and illuminate it. However, this
sunlight must first pass deep through the Earth's atmosphere which filters out most of the
blue colored light. The remaining light is a deep red or orange in color and is much
dimmer than pure white sunlight. Earth's atmosphere also bends or refracts some of this
light so that a small fraction of it can reach and illuminate the Moon.
The total phase of a lunar eclipse is so interesting and beautiful precisely because of the
filtering and refracting effect of Earth's atmosphere. If the Earth had no atmosphere, then
the Moon would be completely black during a total eclipse. Instead, the Moon can take
on a range of colors from dark brown and red to bright orange and yellow. The exact
appearance depends on how much dust and clouds are present in Earth's atmosphere.
Total eclipses tend to be very dark after major volcanic eruptions since these events
dump large amounts of volcanic ash into Earth's atmosphere. During the total lunar
eclipse of December 1992, dust from Mount Pinatubo rendered the Moon nearly
invisible.
All total eclipses start with a penumbral followed by a partial eclipse, and end with a
partial followed by a penumbral eclipse (the total eclipse is sandwiched in the middle).
The penumbral phases of the eclipse are quite difficult to see, even with a telescope.
However, partial and total eclipses are easy to observe, even with the naked eye.

Rotations
A rotation turns all the points in the plane around one point, which is called the center of
rotation. A rotation does not change the sense of figures in the plane. The center of
rotation is the only invariant point (point that maps onto itself) under a rotation. A
rotation of 180 degrees is called a half turn. A rotation of 90 degrees is called a quarter
turn.
Convex, Concave and Plane Mirrors

Purpose: To determine the focal length of a concave (convex) mirror by two different
methods. To determine the relationship between the angle of rotation of a plane mirror
and the angle of deflection of a ray.
Apparatus: Convex, concave and plane mirrors, protractors, projection lamps.
Procedure:
A. Lay a concave (convex) mirror horizontally on the table. The reflected light from
light rays parallel to the principle axis will converge on (or diverge from) a common
point, called the focus. Using the optics lamp, trace out at least three reflected rays (from
incident parallel rays) to determine the focus. Measure the focal length for both the
concave mirror and the convex mirror.
B. Place a pin object (angled into cardboard so that the colored tip points clearly to the
right) in front of a concave mirror. View the image when the object distance is 5 cm from
the center of the mirror segment. Be sure to view the image at a far distance (at least the
width of the table). Record information about the image: erect or inverted, magnified or
demagnified. Repeat for a pin object placed at 10 cm and increments of 1 cm from 10 cm
until the object distance is 20 cm. Estimate the focal length, noting that the image should
become inverted when the object is placed beyond the focal length.

C. Observe the image formed by a convex mirror. Is it erect or inverted?
D. Place the plane mirror near the center of a sheet of paper. (The mirror may be propped
up by some means or a holder used if available.) Draw a line along the silvered side of
the mirror. Then lay a pin about 10 cm in front of the mirror and parallel to its length
(Fig. 1). Mark the locations of the ends of the pin with a pencil. Stick a reference pin R
in the board to one side of the object pin and near the end of the paper, as illustrated in
the figure, and mark its location. Placing another pin nearer the mirror, align this pin and
the reference pin with the head of the pin image in the mirror. Mark the position of the
pin and label with an "H". Then move this pin and align it and the reference pin with the
"tail" of the image pin. Mark this location and label with a "T". Repeat this procedure on
the opposite side of the object pin with another reference pin.
E. Remove the equipment from the paper and draw straight lines from the reference
points through each of the H and T locations and the mirror line. The H lines and the T
lines will intersect and define locations of the head and tail of the pin image,
respectively. Draw a line between the line intersections (the length of the pin image).
Measure the length li of this line and the length lo of the object pin and record. Also,
measure the object distance do and the image distance di from the mirror line and record.
F. Place the plane mirror near the center of a sheet of paper (as described above) and
draw a line along the length of the silvered side of the mirror. Measure so as to find the
center of the line and mark. Stick two pins (A and B) in the board to one side and in front
of the mirror as in Fig. 2. Viewing the aligned images of these pins, place two more pins
(C and D) in alignment. Label the locations of these pins.
G. Leaving pins A and B in place, rotate the mirror a small but measurable angle 
(approximately 10 to 15C) about its center point and draw a line along the silvered side
of the mirror. Align two pins (E an F) with the aligned images of A and B and mark and
label the locations of E and F.

H. Remove the equipment from the paper and draw the incident ray diagram and the two
reflected rays. Measure the angle of rotation of the mirror, and the angle of deflection,
between the two reflected rays and record.
Analysis:
1. For the plane mirror, the object distance do should equal the image distance di.
Compute the percent difference of the measured quantities di and do. Also, the object
length lo should equal the image length li. Compute the percent difference of the
measures quantities lo and li.
2. From Procedure part H, double  and compute the percent difference between 2 and
. Make a conclusion about the relationship of the angle of rotation of a mirror and the
angle of deflection of a ray.
Questions:
1. From Procedure part A, which mirror, convex or concave, is a diverging mirror? Why?
2. Which of the two parts, A or B, is the most accurate for determining the focal length
of the concave mirror? Use the accepted value of f = 15 cm.
3. For a concave mirror, when the object distance is shorter that the focal length is the
image real or virtual? When the object distance is greater than the focal length, is the
image real or virtual?
4. In Procedure part B, why must the image be viewed from a far distance? Hint: Where
does the image form when the object is near the focal point?
Theory:

Fig. 1 Arrangement for experimental procedure for a plane mirror.
Fig. 2 The experimental procedure for rotation of a mirror.
Total Internal Reflection
When light is incident upon a medium of lesser index of refraction, the ray is bent away
from the normal, so the exit angle is greater than the incident angle. Such reflection is
commonly called "internal reflection". The exit angle will then approach 90° for some
critical incident angle c and for incident angles greater than the critical angle there will
be total internal reflection.

The critical angle can be calculated from Snell's law by setting the refraction angle equal
to 90°. Total internal reflection is important in fiber optics and is employed in polarizing
prisms.
For any angle of incidence less than the critical angle, part of the incident light will be
transmitted and part will be reflected. The normal incidence reflection coefficient can be
calculated from the indices of refraction. For non-normal incidence, the transmission and
reflection coefficients can be calculated from the Fresnel equations.
Total internal reflection

An effect that combines both refraction and reflection is total internal reflection.
Consider light coming from a dense medium like water into a less dense medium like air.
When the light coming from the water strikes the surface, part will be reflected and part
will be refracted. Measured with respect to the normal line perpendicular to the surface,

the reflected light comes off at an angle equal to that at which it entered at, while that for
the refracted light is larger than the incident angle. In fact the greater the incident angle,
the more the refracted light bends away from the normal. Thus, increasing the angle of
incidence from path ``1'' to ``2'' will eventually reach a point where the refracted angle is
90o, at which point the light appears to emerge along the surface between the water and
air. If the angle of incidence is increased further, the refracted light cannot leave the
water. It gets completely reflected. The interesting thing about total internal reflection is
that it really is total. That is 100% of the light gets reflected back into the denser
medium, as long as the angle at which it is incident to the surface is large enough.
Fiber optics uses this property of light to keep light beams focussed without significant
loss.
Figure 10.6: Fiber optic cable
The light enters the glass cable, and as long as the bending is not too sudden, will be
totally internally reflected when it hits the sides, and thus is guided along the cable. This
is used in telephone and cable TV cables to carry the signals. Light as an information
carrier is much faster and more efficient than electrons in an electric current. Also, since
light rays don't interact with each other (whereas electrons interact via their electric
charge), it is possible to pack a large number of different light signals into the same fibre
optics cable without distortion. You are probably most familiar with fibre optics cables in
novelty items consisting of thin, multi-coloured strands of glass which carry light beams
Mirages: A Primer
Parched faces look out across the desert sand as two men stagger and crawl toward what
they believed to be a thirst-relieving pool of water. As they near, however, the waters

disappear; the vision merely a mirage. How many times have such scenes been played
out in early motion pictures? The characters may have varied -- cowboys in Death Valley
of the US Southwest desert or French Legionnaires in North Africa's Sahara Desert -- but
the symbolism was the same: men dying of thirst chasing an image that only existed in
their minds.
My standard English dictionary defines mirage as: "an optical illusion due to
atmospheric conditions by which reflected images of distant objects are seen." The word
finds its origins in the French verb se mirer: to be reflected. I'll excuse the lexicographers
on the use of the word reflected in the definition because it does seem more appropriate,
given the word's roots, even if incorrect. Actually they were only two letters off, the
proper technical term should be refracted. (I checked two other English dictionaries, and
they both used the term reflected in their definition. And be careful, a number of popular
books on weather have fallen into the same trap when describing mirages.)
A correct definition can be found in the American Meteorological Society's Glossary of

Weather and Climate. It states:
"Mirage: A refraction phenomenon wherein an image of some distant object is made to

appear displaced from its true position because of large vertical density variations near
the surface; the image may appear distorted, inverted, or wavering
The effects of these distortions, displacements, etc. of the image create many of the
optical illusions we see in a mirage. Although the mind may misinterpret the mirage
image it receives from our eyes, the image is no figure of the imagination -- it can be
photographed. But because of the strong illusional nature of mirages, they have gained

an air of magic about them, dwelling in story and legend along with gods and demons,
fairies and magicians.
The desert mirage of adventure stories and movies and song (for example, the
country/western classic Cool Water) is more commonly seen today as a highway mirage,
those apparent pools of water lying across the pavement, which always disappear before
we can reach them. Both of these forms are technically known as the inferior mirage, a
form most often seen over a land surface. The other main type of mirage is the superior
mirage, most often viewed around or over large bodies of water or snow/ice fields.
The terms inferior and superior are not commentaries on their quality of the mirage's
appearance but refer to the perceived position of the image relative to the actual location
of the object. Inferior mirages appear below the actual object's true location. Superior
mirages are seen above the actual location.
Mirages form when light rays emitted from a source or reflected off an object are bent as
the path of the light ray crosses air layers of different densities. The technical term for
this bending is refraction. You can perform a simple experiment to illustrate refraction.
Place a long pencil or straw in a clear glass filled with water, then look at the pencil in
the glass from the side. It will appear to have a bend or kink where it enters the water.
The degree of bend defines the medium's index of refractivity and depends on the
medium's density. (Refraction also differs for the various colours (wavelengths) of the
visible spectrum and is part of the process causing many atmospheric optical phenomena
including rainbows.)
The index of refraction for air varies with the density of the air. Air density is strongly
dependent on its pressure, temperature and water vapour content. Air density is
proportional to its pressure (density increases as pressure increases) and inversely
proportional to its temperature (density decreases as temperature increases) and to
moisture content (decreasing as water vapour content increases).

One form of superior mirage is so common that most scientists do not even recognize it
as a mirage. Have you ever seen the sun lying right on the horizon? And what was the
actual position of the sun at that time? Yes, this is a trick question.
The fact is, when you see the full solar disk directly on the horizon, all or part of it is
below the horizon, but the bending (refracting) of the solar rays by the atmosphere gives
us the optical illusion that the sun is actually on the horizon. This happens at both sunrise
and sunset each day and adds over 4 extra minutes of daylight to each day. As a result,
daylight is a little longer than the night period on the date of the Equinox by about 4
minutes. [see Equal Night and Day?]
This effect on the setting/rising sun's position is mostly due to the increase in air density
due to increasing atmospheric pressure as the solar rays come from the near vacuum of
outer space toward the Earth's surface. (Similar effects also alter the observed position of
the moon, planets and stars.) When the sun or moon are within a few degrees of the
horizon, they can also lose their circular shape and appear flattened because of
atmospheric refraction. There are several other refractive effects that are also of interest
to sky watcher but we must leave them for a later date. In this piece I want to focus on
refractive phenomena seen in the lower atmosphere under conditions of varying
temperature with height which produce mirages.
The two main mirage types, superior and inferior, are caused by opposite patterns of
temperature change with height -- what meteorologists call the vertical temperature
gradient. They express this gradient by the number of degrees of temperature change per
height increment (usually Celsius degrees per metre or per 100 metres, but in older
books you will sometimes see it still expressed as either Fahrenheit degrees per foot or
per 100 feet).

Meteorologists consider temperature decreasing with height a "normal" state of the

atmosphere -- often using the term lapse rate for the condition. When temperature
increases with height, the condition is termed a temperature inversion (or just inversion).
Inferior mirages occur under strong lapse rate conditions, and superior mirages under
inversion conditions.
While the atmosphere tries to establish a uniform lapse rate of 0.98 C° per 100m (0.54
F° per 100 ft) throughout the lower atmosphere, this situation rarely occurs for long as
heating and cooling establishes differing temperature gradients from the surface upward.
At times there may be several layers with distinct temperature gradients including
inversion layers below or between layers with lapse conditions.
Such complex temperature layer can cause some rather interesting optical situations such
as the Novaya Zemlya mirage which can "carry" light rays hundreds of kilometres over
the horizon. We will delve more deeply into the variations possible with mirages in
separate articles on this website. For the remainder of this piece, we will focus on
describing the general conditions that form the inferior and superior mirages.
In the lower layers of the atmosphere where most mirages are seen, the air refractivity is
only weakly dependent upon changes in air pressure but strongly on temperature
gradient changes.
Light reaching our eye, after traveling through a region of the atmosphere where the
temperature gradient is constant, follows a curved (parabolic) path. (See diagrams below.
In them, the dark lines indicate the actual light ray path and the white dashed lines the
path our mind thinks it sees.) The degree of curvature is proportional to the temperature
gradient along that path -- the stronger the temperature gradient through which the light

passes, the greater the bend. When the temperature gradient is not constant but changes
with height, the curvature of the light path can increase more rapidly at some heights
than others, thus producing interesting effects such as object distortions and multiple
images.
The light ray always bends toward the colder (and thus denser) air, so that the colder air
is on the inside of the curvature. The image we see is always displaced in the direction of
the warmer air. Therefore, if we have a temperature gradient with warmer air nearest the
surface, the image will be displaced downward toward the surface -- forming an inferior
mirage. When the air is colder near the surface (an inversion condition), the image wis
displaced upward, forming a superior mirage.
Mirages are confined to small viewing angles even when they appear large, about half a
degree in width -- the size of the solar disk -- and most portray objects located from half
a kilometre to about five kilometres (about a quarter mile to three miles). Under strong
inversion conditions, however, objects hundreds of kilometres away can be seen,
including those located beyond the normal viewing horizon. (For more details on seeing
beyond the horizon, see The Arctic Mirage: Aid to Discovery.)
The illusion part of the mirage generally come from our mind's interpretation of what the
eye sees. When our eye sees a light ray coming from an object, our mind interprets the
ray path as a perfectly straight line from the object to our eye. Only when we recognize
that we are seeing an illusion can we make some mental corrections to the scene. For
example, we know the pencil in the glass is not bent even if our eyes try to tell us so. But
recall what happens initially when we are surprised by such a scene. We must process
more information to make true sense of the scene before our eyes. Thus when we see the
inferior mirage's pool of water on the ground, we expect that the "water" will actually
on. the ground

When we view an inferior-mirage "water pool" on a road or desert or other hot surface,
what we are actually seeing is the image of the bluish sky being strongly refracted by the
hot air near the surface so that it appears to our mind to be water lying on the surface.
The appearance of the inferior mirage always indicates that the surface air is much
warmer than the air above it due to the strong heating of the surface by the sun or some
other heat source. This temperature structure -- very hot below and cool above -- causes
light rays passing through it to be bent upward.
The superior mirage occurs under reverse atmospheric conditions from the inferior
mirage. For it to be seen, the air close to the surface must be much colder than the air
above it. This condition is common over snow, ice and cold water surfaces. When very
cold air lies below warm air, light rays are bent downward toward the surface, thus
tricking our eyes into thinking an object is located higher or is taller in appearance than
it actually is

.
The superior mirage can also make objects appear to be floating in the air or cause
objects actually located below the horizon to appear above it (remember the setting-sun
example), a condition called looming. The superior mirage can also cause objects appear
to be taller than they actually are, called towering, or shorter, a condition termed
stooping.
Occasionally, we can see a variation on the inferior mirage on a smaller scale over a
vertical surface. In this situation, the strong heating of a vertical surface by the sun or an
internal heat source (such as a motor cover) can develop a strong temperature gradient
extending laterally outward from the surface. This condition can form a lateral mirage.
A lateral mirage will appear as an apparent reflection of a nearby image and form just
over the wall or rock face.
This "inferior" lateral mirage results from a strong temperature gradient next to the wall
similar to the condition of an inferior mirage turned sideways. David K. Lynch and
William Livingston in their book Color and Light in Nature suggest that there is no
reason a "superior" lateral mirage could not form over a cold wall surrounded by warm
air, but they had never seen one.
Types of Mirages
Two-image mirages (these are the classical types):
1. Inferior mirage
So called because the inverted image is below the erect one. This is the familiar hot-road
mirage seen every sunny day on smooth paving. Caused by the thin layer of hot air
below eye level, at the surface.

2. Superior mirage
Here the inverted image is above the erect one. Caused by a layer of hot air not far
overhead (a "thermal inversion".) Like the inferior mirage, this looks much like a simple
reflection. Astronomical objects cannot appear in a simple superior mirage.
Three-image mirages:
An inverted image lies between two erect ones. The top image is often strongly
compressed. These purely refractive phenomena are also caused by inversion layers;
they are of at least two kinds:
3. The “mock mirage”

Caused by looking down into an inversion below eye level, and then (thanks to the
curvature of the Earth) out through it again beyond the horizon. The miraged objects
may be about the same height above sea level as the eye, or may be considerably higher.
4. Wegener's “late mirage”

Caused by looking up through an inversion above the observer. The miraged objects are
always higher than eye level (e.g., distant mountains; astronomical objects). A true
superior mirage of objects below the inversion may also be present, if the inversion is
strong enough.
5. The Fata Morgana:

Is the general name for these; but the phenomena are so varied that two or more mirage
types may really be involved. The image is marked by repeated vertical and horizontal
features, probably due to many repeated alternations of erect and inverted images of
some object. Often the mirage shows considerable internal motion, producing an illusion
that people or animals appear in the scene. Almost certainly, strong inversion layers are
responsible; but the detailed explanation of these rare displays has yet to be found.
Other uses of the term “mirage”:

Sometimes, you will come across the term “lateral mirage,” which is used in two senses.
The first is a supposed sideways displacement of a miraged image, often by many
degrees along the horizon. Such displacements are physically impossible; refraction is
almost entirely in the vertical direction. Reports of large lateral displacements are the
result of mis-identifications. (For example, a distant mountain that is normally hidden by
closer ones may become visible by looming. As many mountains in any small region
have similar shapes, the unfamiliar object is often incorrectly “identified” as a strongly
displaced image of some familiar feature of the normal landscape.)
A more legitimate use of the term refers to mirages seen on sunlit walls. However, these
mirages are simply the familiar inferior (hot-surface) mirage, turned through 90°. In this
case, to avoid confusion with the erroneous reports mentioned above, I would prefer the
term “mural mirage” to be used instead of “lateral mirage.”
Of course, “mirage” is often used metaphorically, to denote a false hope like that of the
thirsty traveler in the desert, who imagines that water lies in the distance on seeing an
inferior mirage. If you search for the word “mirage” on the Web, most of the pages you
turn up contain such rhetorical uses of the term, and do not refer to real mirages.
An optical fiber
Is a flexible filament of very clear glass and is capable of carrying information in the
form of light. This filament of glass is a little thicker than a human hair.
Core and Cladding
The two basic elements of optical fiber are its core and cladding. The core, or the axial
part of the optical fiber, is the light transmission area of the fiber. The cladding is the
layer completely surrounding the core. The difference in refractive index between the
core and cladding is less than 0.5%. The refractive index of the core is higher than that

of the cladding, so that light in the core strikes the interface with the cladding at a
bouncing angle and is trapped in the core by total internal reflection
Multimode vs. Single-mode
A mode is a defined path in which light travels. A light signal can propagate through the
core of the optical fiber on a single path (single-mode fiber) or on many paths
(multimode fiber). The mode in which light travels depends on geometry, the index
profile of the fiber, and the wavelength of the light.
Single-mode fiber has the advantage of high information-carrying capacity, low

attenuation and low fiber cost, but multimode has the advantage of low connection and
electronics cost that may lead to lower system cost
The Prism Prison

Pure light passing through a prism releases a spectrum of color. If you were bathed in red
you would see the prism as a giver of red. If you were bathed in blue, you would see the
prism as a giver of blue, and so on. The blue would vigorously argue that the prism,
although looking white or clear only contains blue, as that is all blue can see and
experience: likewise with red, but of course now red. Blue and red may tolerate the
others vision, but to each the others vision is colored. Blue or red each experiences their
vision as correct. By being correct the others vision becomes faulty. Each may be
tolerant of the others vision, but that other is seen and taken as a misguided vision or a
lesser knowing brother of the prism.
Both blue and red present their experience as the way, and each is correct within the
color spectrum that they live within. The more zealous of the blue and red are compelled
to convert all colors to their color. As their color is the color of the prism, and all other
colors arise from defective vision: tolerable but deficient of true substance. The color
spectrum of the converter initiates tunnel vision, as only its color is seen as true

substance. All other colors are seen as emanating from a lesser prism, and potential
convert's to the one true color.
When blue and red confront each other the only common ground they have is that they
are both colors. Each will represent their color, stained by belief that the color they live
within is the only color. A color is simply one color of the prism, but the belief-stain
creates the potential to make one color superior to the other. The color by itself contains
no conflict, as all colors exist within the Prism with no hostility to each other. It is the
belief-stain that initiates problems that can grow into antagonistic confrontation. The
animosity is stain vs. stain, but the color itself is malevolence free. Stain vs. stain
contains the potential to become violent and destructive, as the stains are in opposition,
even though the colors remain impartial and unaffected. The potential violence can be
deathly confrontational, and many stains can be removed, but the color of each remains
unchanged and unmoved.
The Pure Light of existence becomes visible through the prism of individual existence.
The color each individual lives within, if stain free, will ultimately lead that individual to
the clarity of the color created by their prism, and there discover all the colors existing in
perfect and undisturbed harmony and peace. The Prism is the Mother and the Light is the
Father for the rainbow of existing colors.
Each color points to the prism as proof for the existence of their color. But the belief-
stain prevents each color from recognizing the purity and clarity of the prism light that
originates all the colors. To take one color of a rainbow, and claim that it is the rainbow,
creates the belief-stain. However, the rainbow exists with or without the stain, as its
colors are the natural exposure of pure Light in multiple forms. When the light of the
prism is experienced within absolute clarity, all the colors simply become a spectrum of
harmonious and interconnected colors. Each color blends into the next, with no defined
border. The homogenous blending of one color to another is by altering its hue of each to
meet the next color that has done likewise. They meet in a borderless region where the
two colors are of the same color: the resolution of duality simply achieved by flexibility

of expandable borders. The light of the prism contains its own duality resolution: it's
innate and effortless.
You are both the physical prism (body) and the non-physical light (consciousness)
experienced as a color (individual identity) in the world we now inhabit. The Light
separates through the Prism creating prisms of light - your individuality - which is a
refraction of the Light into multiple colors. The refraction of the Prism is a visible
reflection of the Light that is experienced as both individual light and color. You take the
color as what and who you are, as it can be seen. When the individual prism disintegrates
the color disappears, but your light has no where to go. It remains as it always has,
without or without a prism to see itself as a reflected color. Your light is your reality,
your permanence and your innate clarity.
Experientially recognizing that your light - a reflected refraction of the Light - contains
no darkness unveils your reality. It is Light and nothing but Light, as your light is light
and nothing but light. The darkness is the stain upon the transitory but visible color. As
the color is transitory so is the stain. When the prism goes the color goes, and so does the
stain. But individual light remains, as it has created a new identity for itself, if you have
made that light your life.
In short, there are no hell realms. They are a created fiction to control by fear. Your light
contains no darkness other that what your transitory identity creates. And that created
darkness goes when the transitory goes. The natural purity and brilliant clarity of your
light is your birthright. The way home is always within you. The pathway leading home,
back to your natural and innate primal brilliant clarity has always existed within
yourself. It simply requires unveiling.
It is the belief-stain that creates prism-prisons by attempting to reveal the Light of the
Prism as the color your light exists within. When you bypass, let go or stand aside from
the color, remaining stain free, your light in all its clarity is unveiled, becoming visible
as Light. Color and stain free, the Light shines though you as pure light, which is then,

color free. Or, the Light creates such intense brilliance to the color that it obscures the
color, making only brilliant clarity visible, which is pure Light. And yet, the color of
others, stained or blinded by the darkness of its own color, still can only see color where
no color exists.
The original Prism never alters, only how each individual light perceives the color they
live within alters, which is the originating force for the stain. Without the Prism, nothing
but Absolute Light would be. And even that Light cannot see itself unless reflected.
Individual prisms expose the Light of the original reflection, making it visible as a
refracted reflection, which is you. The color that is seen is the visible reality of your
individual existence, which is clear Light made visible by a refracted reflection through
individual prisms of light.
To suggest live in the Light without knowing how to bypass or stand aside from the color
of the prism creates the same myopic vision as a stain does. To allow the light of your
existence to exist within its natural brilliance is effortless. But effort it required to
achieve that effortless, revealing how the color is constructed and how it functions. If
you do not know how the color is constructed you will not know how the stain appears.
You will have a trying, if not impossible task to experience the clarity of the Light, while
still perceiving yourself as a color, as you tend to color the Light with the color you live
within.
Experientially knowing how the light of the prism is constructed will reveal how it
functions. Or, experientially seeing how it functions exposes its primal construction. In
that primacy the Light is experienced as both refracted and whole. The refraction is seen
as color, whilst the wholeness is your light as a reflection of the Light. You have it all
from the very beginning. The color is the veil and the stain the darkness that appears to
dim your light. But your light is always full on, just veiled.
The prism with ease can become a prison, as the color can enthrall, captivate and entrap.
Removing the stain of your color or not allowing the color to blind you allows the colors
natural purity and brilliance to illuminate the prism. That experiential purity and

brilliance points directly at the clarity from which the color arose, leading you straight
into the essence of the Light. Trusting the purity of the color removes the stain that has
been sustained by fear. Trusting your own trust unveils the Light of the prism. Both
trusts initially take effort, but when the clarity of the Light is revealed, trust becomes
effortless. The Light is trust, as it trusts you with a color, and when your trust is in
absolute harmonic resonate alignment with that trust, you are living within the light-
Light of the prism-Prism, as pure light-Light and not color.
If others only see color where pure light exists, it is their blindness, which can never in
any way diminish the light. Trusting the Light of your existence creates color blindness,
as the Lights brilliance eliminates your ability to discriminate the rainbow colors
emanating from the Prism, as you only see Light wherever you turn.
Only seeing the rainbow emanating from the Prism veils your vision of the Light, as you
can only sees colors. Trusting your innate colorblindness reveals the Light within the
color, as the Light creates, sustains and is within all the rainbow colors of the prism. It is
only the color that wants to color the Light as itself. The Light is color free. To know and
experience the Light as it really is, simply allow it to shine through you, color free. The
Lights natural brilliance then becomes your natural brilliance, and that illumination
blinds you to the rainbow colors. Or, you remain as the light of the prism facing the
Light, and so cannot see colors, only Light. Both processes are effortless.
PRISM
If light passes through a prism, a transparent object with flat, polished surfaces at angles
to one another, the exit ray is no longer parallel to the incident ray. Because the
refractive index of a substance varies for the different wavelengths, a prism can spread
out the various wavelengths of light contained in an incident beam and form a spectrum.
In Fig. 5, the angle CBD between the path of the incident ray and the path of the
emergent ray is the angle of deviation. If the angle the incident ray makes with the
normal is equal to the angle made by the emergent ray, the deviation is at a minimum.
The refractive index of the prism can be calculated by measuring the angle of minimum
deviation and the angle between the faces of the prism.

Total reflection
Let the wall (out of glass) which light encounters as it spreads be perfectly transparent.
This means: The light, as it arrives at the surface of the wall, passes through the
interface between the glass and air without being weakened. However, this never
happens, since one part of the light is always reflected, even at the surface of so
transparent substances as water and glass, into the substance from which it has arrived.
The surface of water reflects, so does every pane of glass like the mirrors of shop
windows. However, the reflected light is always only a small part of the incident light,
while it becomes the larger, the more inclined the rays are to the surface they encounter.
Such reflected images have therefore always very weak light. Their formation is always
accompanied by refraction, which is the principal part of the incident light passing
through the interface.
Note: During this process, the light comes from air and enters water or glass. However,
it is quite different when it comes from water or glass and enters air. When a ray of light
comes from air and enters water or glass, the angle of refraction is always smaller than
the angle of incidence; however large is the angle of incidence (Fig. 661, left image,
between 0º and 90º), there always is added an angle of refraction. However, it is different
when the light comes from water or glass and enters air. The right image in Fig. 661
explains: The angle of incidence r in water is then largest when the angle of refraction in
air is a right one, because this is the largest angle, which the ray can form with the
perpendicular NN in the semi-circle of air. This angle r is the limiting angle. The rays in
Fig. 662, which form a still larger angle with the normal, for example, g, cannot at all
exit into the air, they are reflected into the water, that is, they are reflected at the
interface of water and air. This process is called total reflection, because all rays are

reflected, as you conclude from the fact that the reflected light has the same strength as
the incident light. In brief: If light passes from a more strongly refracting substance in to
a more weakly refracting one, that is, into a substance, in which the ray is refracted away
from the perpendicular and exceeds the angle of incidence by a certain value, the light is
reflected altogether back into the more strongly refracting substance. - How large is the
limiting angle? If in air the angle of incidence i is like in Fig. 661 (left image) a right
one, then sin 90º/sinr = n, and since sin 90º = 1, sin r = 1/n
The angle r is the limiting angle: Imagine the sequence in which the ray passes through
the substances in the reverse order (Fig. 661 right image), you see then that r is the angle
of incidence, to which is added an angle of refraction: The emitting ray touches the
surface. For water, n = 1.33, whence the limiting angle is that angle the sin of which is
3/4, that is, the angle 48º 36' 25; for one kind of glass (light crown glass ) with n = 1.50,
the limiting angle is 41º 48' 37|.
You can observe total reflection at a water surface which you view - as in an aquarium -
from down below, but also by looking at it in a glass of water from below. You cannot
see through the water surface and get the impression that it is a perfect mirror.
The total reflection of light explains how certain objects which as a rule are shiny and at
the same time transparent, become under certain conditions opque and dull, for example,
powdered glass and ice, snow, foam (irrespectively whether it is made out of coloured or
colourless substances). Powdered glass is, strictly speaking, a mixture of glass and air. It
is dull for the same reason as, for example, polished silver is, when it is powdered: It
turns its reflecting faces in all possible directions, that is, throws light in all directions,
that is, it is diffusive. It is opaque due to the total reflection of light in the mixture of
glass and air; light cannot pass through the mixture. Powdered glass becomes
transparent, if you cover it with oil, which has approximately the same refractive ratio as
glass

Binoculars
How binoculars work
Essentially, binoculars are just two telescopes mounted side by side, one for each eye. To
understand binoculars, you need to understand how a telescope works. Here's an easy
demonstration that you can try yourself. All you need are two ordinary magnifying
glasses and a piece of tracing paper. Do this once, and you will understand forever how
binoculars work.
Hold the tracing paper on the opposite side of the magnifying glass from a bright object,
such as a light bulb. Move the paper back and forth. At a certain distance, an upside-
down-and-backwards image of the light bulb will form on the paper.
You can enlarge this image by examining it through another magnifying glass, as shown
in the illustration. You may be surprised to find that if you slide the tracing paper away,
the image will remain, only brighter and clearer. You have just made a working
telescope.The magnifying glass nearest the object is called the objective lens; the one
nearest your eye, the eyepiece. The objective lens and the eyepiece are two elements in
all binoculars. Binoculars also have a third element, the erecting prisms.

In the telescope we just built, everything is upside down and backwards. That would be
OK for looking at stars, but for watching birds or following the action at a football game
we require a right-side-up picture. A terrestrial telescope has to flip the image, and that's
what prisms do.
A prism is a solid piece of glass that functions as a mirror, but without a mirror's
reflective backing. Light rays that have entered a prism cannot get out if they strike a
surface at too great an angle. Instead, they reflect back, as if from a perfect mirror.
In the mid 19th century, an Italian named Porro designed a telescope with two prisms set
at right angles to each other between the objective lens and the eyepiece. This
arrangement not only erected and reversed the image, but also folded the light path,
resulting in a shorter, more manageable instrument. In 1894, the Zeiss Optical Works
created the first "Hunting Glasses," incorporating the Porro prism design, and modern
prismatic binoculars were born.
All binoculars still have these three parts. An objective lens focuses an upside-down
image. A set of prisms turns the image right side up. And an eyepiece magnifies it.
Though modern eyepieces and objective lenses are each comprised of multiple elements,
their basic functions remain unchanged.
Today you can buy binoculars that are made with roof prisms or Porro prisms. There are
advantages to each.

Diffraction of Light
In his 1704 treatise on the theory of optical phenomena (Opticks), Sir Isaac Newton
wrote that "Light is never known to follow crooked passages nor to bend into the
shadow". He explained this observation by describing how particles of light always
travel in straight lines, and how objects positioned within the path of light particles
would cast a shadow because the particles could not spread out behind the object.
On a large scale, this hypothesis is supported by the seemingly sharp edges of shadows
cast by rays from the sun. However, on a much smaller scale, when light waves pass near
a barrier, they tend to bend around that barrier and spread at oblique angles. This
phenomenon is known as diffraction of the light, and occurs when a light wave passes
very close to the edge of an object or through a tiny opening such as a slit or aperture.
The light that passes through the opening is partially redirected due to an interaction
with the edges. An example of light diffraction is presented in Figure 1 for coherent red
laser light passing through a very tiny line grating composed of a series of bars on a
glass microscope slide. The bars diffract the laser light into periodic beams of light that
can be visualized in the figure. Diffraction is a phenomenon similar to dispersion but is
not related to a variation in the wavelength of light.
Bright bands that are often seen inside the edges of geometric shadows are the result of
diffraction. When light waves originating from a distant point of light strike an opaque
object, they tend to bend around the edges, curving both into the shadow and back
through the path of other light waves from the same source. The waves that bend around
behind the object create a bright line where the shadow would ordinarily begin, but

waves that bounce back into the path of the light overlap waves from the source, creating
an interference pattern of light and dark bands around the edge of the object (see
Figure2) Diffraction is often explained in terms of the Huygens principle, which states
that each point on a wavefront can be considered as a source of a new wave.
Depending on the circumstances that give rise to the phenomenon, diffraction can be
perceived in a variety of different ways. Scientists have cleverly utilized diffraction of
neutrons and X-rays to elucidate the arrangement of atoms in small ionic crystals,
molecules, and even such large macromolecular assemblies as proteins and nucleic
acids. Electron diffraction is often employed to examine periodic features of viruses,
membranes, and other biological organisms as well as synthetic and naturally occurring
materials. No lens exists that will focus neutrons and x-rays into an image, so
investigators must reconstruct images of molecules and proteins from the diffraction
patterns using sophisticated mathematical analysis. Fortunately, magnetic lenses can
focus diffracted electrons in the electron microscope, and glass lenses are very useful for
focusing diffracted light to form an optical image that can be easily viewed.
A very simple demonstration of light diffraction can be conducted by holding one hand
in front of a strong light source and slowly bringing two fingers close together while
observing the light transmitted between them. As the fingers approach one another and
come very close together (almost touching), one can begin to see a series of dark lines
parallel to the fingers. The parallel dark lines together with the bright areas between
them are actually diffraction patterns. This effect is clearly demonstrated in Figure 2,
for diffraction rings that appear surrounding the sharp edges of a razor blade when it is
illuminated with intense blue light from a laser source.
Another simple, but very common, example of diffraction occurs when light is scattered
or bent by small particles having physical dimensions in the same order of magnitude as
the wavelength of light. A good illustration is the spreading of automobile headlight
beams by fog or fine dust particles. The amount of scattering and the angles taken by the
redirected light beams are dependent upon the size and density of the particles causing

the diffraction. Light scattering, a form of diffraction, also underlies the blue color of the
sky and the often beautifully colored sunrises and sunsets that can be observed on the
horizon. If the Earth were devoid of an atmosphere (lacking air, water, dust, and debris),
the sky would appear black. When light from the sun passes through the Earth's
atmosphere, gas molecules having varying densities, due to temperature fluctuations and
the amount of water vapor present, will scatter the light. The shortest wavelengths (violet
and blue) are scattered to the greatest extent, rendering the sky a rich, deep blue color.
When there is a considerable amount of dust or moisture in the air, longer (primarily red)
wavelengths also become scattered along with the blue wavelengths, causing the blue
sky to become whiter in color.
When the sun is high (around noon) in a clear dry atmosphere, most of the visible light
passing through the atmosphere is not scattered to a significant degree, and the sun
appears almost white on a deep blue background. As the sun begins to set, the light
waves must pass through increasing amounts of atmosphere, usually with larger
quantities of suspended dust and moisture. Under these circumstances, longer
wavelengths of light become scattered and other colors start to dominate the color of the
sun, which ranges from yellow to orange, finally turning red just before it drops below
the horizon.

We can often observe pastel shades of blue, pink, purple, and green in clouds, which are
generated by a combination of effects when light is refracted and diffracted from water
droplets in the clouds. The amount of diffraction depends on the wavelength of light,
with shorter wavelengths being diffracted at a greater angle than longer ones (in effect,
blue and violet light are diffracted at a larger angle than is red light). The terms
diffraction and scattering are often used interchangeably and are considered to be
almost synonymous in many cases. Diffraction describes a specialized case of light
scattering in which an object with regularly repeating features (such as a periodic object
or a diffraction grating) produces an orderly diffraction pattern. In the real world, most
objects are very complex in shape and should be considered to be composed of many
individual diffraction features that can collectively produce a random "scattering" of
light.
In the microscope, scattering or diffraction of light can occur at the specimen plane due
to interaction of the light with small particles or features, and again at the margins of the
objective front lens or at the edges of a circular aperture within or at the rear of the
objective. It is this diffraction, or spreading of light, that makes it possible to observe
magnified images of specimens in the microscope, however it is also diffraction that
limits the size of objects that can be resolved. If light passes through a specimen and is
not absorbed or diffracted, the specimen will not be visible when viewed through the
eyepieces. The manner in which an image is formed in the microscope depends on the
diffraction of light into divergent waves, followed by their subsequent recombination
into a magnified image through constructive and destructive interference.
When we view a specimen, whether directly or with a microscope, telescope, or other

optical instrument, the image we see is composed of a myriad of overlapping points of
light emanating from the plane of the specimen. Therefore, the appearance and integrity
of the image from a single point of light holds a significant amount of importance with
regards to formation of the overall image. Because the image-forming light rays are
diffracted, a single point of light is never really seen as a point in the microscope, but
rather as a diffraction pattern containing a central disk or spot of light having a finite

diameter and encircled by a fading series of rings. As a result, the image of a specimen is
never an exact representation of the specimen, and a lower limit is imposed on the
smallest detail in the specimen that can be resolved. The resolving power is the ability of
an optical instrument to produce clearly separated images of two adjacent points. Up to
the point at which diffraction causes the resolution to be limited, the quality of the lenses
and mirrors in the instrument, as well as the properties of the surrounding medium
(usually air), determine the final resolution.
Several of the classical and most fundamental experiments that help explain diffraction
of light were first conducted between the late Seventeenth and early Nineteenth
Centuries by Italian scientist Francesco Grimaldi, French scientist Augustin Fresnel,
English physicist Thomas Young, and several other investigators. These experiments
involved propagation of light waves though a very small slit (aperture), and
demonstrated that when light passes through the slit, the physical size of the slit
determines how the slit interacts with the light. If the wavelength of light is much
smaller than the aperture or slit width, a light wave simply travels onward in a straight
line after passing through, as it would if no aperture were present (as presented in
Figure 3).

However, when the wavelength exceeds the size of the slit, diffraction of the light
occurs, causing the formation of a diffraction pattern consisting of a bright central
portion (the primary maximum), bounded on either side by a series of secondary
maxima separated by dark regions (minima; see Figure 4). The maxima and minima are
created by interference of diffracted light waves. Each successive bright band becomes
less intense proceeding outward, away from the central maximum. The width of the
central bright portion, and the spacing of the accompanying sidebands, depend on the
size of the aperture (slit) and the wavelength of the light. This relationship can be
described mathematically and demonstrates that the width of the central maximum
decreases with decreasing wavelength and increasing aperture width, but can never be
reduced to the size of a point light source.
The intensity distribution of light diffracted by the single slit experiment is presented in
Figure 4. It is assumed that both light beams in Figure 3 are composed of coherent,
monochromatic waves emitted from a point source that is far enough away from the slit
for the wavefronts to be considered linear and parallel. Light passing through aperture d
on the right-hand side of the figure has a wavelength larger than the aperture and is
diffracted, with the primary incident light beam landing at point P and the first
secondary maximum occurring at point Q. As shown on the left-hand side of Figure 4,
when the wavelength is much smaller than the aperture width (d), the wave simply
travels through in a straight line, just as it would if it were a particle or no aperture were
present. However, when the wavelength exceeds the size of the aperture, it is diffracted
to produce a central peak containing most of the light intensity with secondary higher-
order maxima and intensity minima governed by the equation:
sin() = n/d
where  is the angle between the central incident propagation direction and the first
minimum of the diffraction pattern, and n indicates the sequential number of the higher-
order maxima. The light intensity is maximum at  = zero degrees and decreases to a
minimum (where the intensity is zero) at angles dictated by the equation above. The
experiment produces a bright central maximum, which is bounded on both sides by

secondary maxima, with the intensity of each succeeding secondary maximum

decreasing as the distance from the center increases. Figure 4 illustrates this point with a
plot of beam intensity versus diffraction radius. Note that the minima occurring between
secondary maxima are positioned in multiples of .
Both the experiment described above, and the demonstration of diffraction using light
passing between the fingers, utilize a narrow slit as the aperture producing a diffraction
pattern. All optical instruments, including microscopes, utilize circular lenses and
apertures, as does the human eye itself. Circular apertures produce similar diffraction
phenomena, although with circular symmetry (instead of linear geometry, as in the case
of slits). Therefore the diffraction pattern of a point source of light, if highly magnified,
is seen to consist of a central bright disk surrounded by a series of diffraction rings (the
secondary maxima and minima). When the lens, such as a microscope objective lens, is
properly focused, the light intensity at the minima between the bright rings in the pattern
is zero. No matter how perfect the lens is, the secondary diffraction maxima cannot be
eliminated nor can the central spot be reduced to a single point of light (unless the lens
could be made to be of infinite diameter).

The central diffraction spot or disk is called an Airy disk, named after Sir George Airy,
who described many aspects of the concept in the Nineteenth Century. The Airy disk
pattern (illustrated in Figure 5) is a direct result of diffraction, and demonstrates the
alteration of points of light that make up an image using an optical instrument such as a
microscope. In a manner that is similar to diffraction by a slit, the size of the central disk
produced by circular lenses is related to the wavelength of the light and the diameter or
aperture angle of the lens. In the case of a camera or telescope lens receiving light from
an object at great (infinite) distance, the aperture angle depends on the focal ratio f/D,
where D is the lens diameter and f is the focal length. The focal ratio is usually referred
to in photography as the f-number of the lens. The aperture angle can be considered to be
the angular diameter of the lens, measured from a reference point at the lens aperture to a
point in the image plane positioned at the focal length (f) away from the lens. The radius
of the diffraction disk (d) is given by the relationship:
d = 1.22 • (f/D)
With objective lenses employed by the microscope, the term numerical aperture (NA) is
used instead of aperture angle. The definition of numerical aperture includes the
refractive index of the medium positioned between the front of the lens and the
microscope specimen slide, and the half angle over which the lens can collect light from
a nearby specimen placed at the focal distance. Using the term n to designate the
refractive index, and  for the half angle aperture, the numerical aperture is defined as :
NA (Numerical Aperture) = n • sin ()
The radius of the diffraction spot (r) for a point of light in the image plane (see Figure 4)
is given by the related expression:
r = 1.22 • /(2NA)
]The Airy disk patterns, along with point spread functions, at three hypothetical
resolutions are presented in Figure 5. The point spread function is a three-dimensional
representation of the diffraction pattern occurring along the optical axis of the
microscope. As the lateral resolution increases, Airy disk size decreases and the

corresponding point spread function narrows. This can be demonstrated in observing the
figure by comparing the Airy disk and point spread function in (a), which display the
lowest resolution, with the corresponding set in (c) that have the highest resolution of the
group. Experimentally, resolution can be increased by decreasing the wavelength of light
utilized to image the specimen (from white light to blue, for example) or by increasing
the numerical aperture of the objective and condenser combination. Under most
circumstances, it is easier and far more practical to select an objective having a higher
numerical aperture in order to increase the resolution of the images produced by the
microscope.
Regardless of whether an image is formed in the microscope or another optical

instrument, the size of a diffracted point of light becomes smaller with decreasing
wavelength and increasing numerical aperture, but always remains a disk that is larger
than the point of light originating from the specimen (or other object) being imaged. In
evaluating the resolution that is possible with a microscope, if the size of the individual
diffraction spot is the limiting factor (rather than lens aberrations or other variables), the
image achieved is said to be diffraction limited. For any optical instrument, therefore,
the light gathering ability is fixed by the aperture angle or numerical aperture, and the
resolution obtained is controlled by varying these values and the wavelength of light
utilized to capture the image in order to obtain the smallest diffraction disk size that is
possible with the instrument. Only when specimen details seen in the image are larger
than this limiting disk size, can conclusions be made about the size, shape, and
arrangement of the features.

LIMIT OF RESOLUTION
The maximum resolving power of a far-field optical microscope is given by the Rayleigh
criterion, which we also refer to as the diffraction limit:
Where is the wavelength of light and is the numerical aperture. Eq. 2.29 is
derived by Fourier-optics theory (see next section): Points in real space appear as Airy
discs in the image plane of a microscope. In this discussion we shall assume an ideal
microscope without any aberrations. The resolution dmin refers to the distance between
two points in real space for which the maximum intensity of one Airy disc falls on the
first minimum of the second (see Fig. 2.9). At this distance, the two points are said to be
well-resolved by the microscope.
Figure 2.9: Imaging of two points at the resolution

limit of a far-field optical microscope. The image is a
superposition of two Airy discs with the maximum of
one Airy disc at the position of the first minimum of
the second.
It should be mentioned that the actual experimental resolution is lower than the
diffraction limit in general. Non-ideal optical elements within the microscope cause
aberrations, and thus reduce resolution. Additionally, resolution can be diminished due to

low contrast and noise. The latter aspect has to be taken into account, e.g., in magneto-
optical microscopy: If polar MOKE is detected, the effect is strongest for normal
incidence and declines with increasing angle of incidence. However, for high lateral
resolution, a large aperture has to be chosen. The magneto-optical contrast is therefore
diminished, affecting both the sensitivity and the experimental resolution in the final
magneto-optical images
In high-resolving microscopes, oil-immersion objectives with apertures are
used, resulting in a maximum resolution of . The Rayleigh criterion is

derived under the assumption of incoherent light. For coherent light, under certain
circumstances, the resolution can slightly be enhanced to . However,

(non-scanning) coherent light microscopy is difficult since the images are often
superimposed with diffraction patterns that originate from small structures in the optical
path. Highest resolution is obtained by scanning confocal microscopy. Here, the image
of an illuminating aperture in the sample plane coincides with the "back-focused" image
of the exit aperture. Although the size of the illumination and detection spots are
diffraction limited, the effective resolution obtained by the combination of the two (as
determined by the width of the product of the two respective Airy discs) is enhanced by
a factor . For recording images of a specimen, the illumination spot is scanned by

translating the sample. In addition to high lateral resolution, confocal microscopy also
permits imaging of planes at varying depths within the specimen.
Limits to Resolution in the Transmission Electron Microscope

It is desirable to understand several of the fundamental principles of light optics in order

to understand the limitations of electron microscopy.
Diffraction. First, it is impossible to achieve absolute focus using any optical system
that uses particles with wave-like properties, because of diffraction and interference.
Diffraction results when a wavefront is impeded by any object, and of course the edge of
the lens area constitutes an object, as does any superimposed aperture. Below is an
example of how diffraction changes the wavefront in the presence of a small aperture.
Notice that this causes the parallel wavefront to emerge from the aperture as a spherical
wavefront.
Airy disc. Second, using even a "perfect" optical system, a point of light cannot be
focused as a perfect dot. Instead, the image when viewed critically consists of a disc
composed of concentric circles with diminishing intensity. This is known as an Airy disc
and is represented below.

Resolution. Notice the primary, secondary and tertiary wavefronts generated by the Airy
disc. (Of course, these continue to emanate at higher orders, but their affect on optical
phenomena diminishes in importance with each higher order.) The resolution is typically
described as the distance between the first order peak and the first order trough
(designated "r" above). Resolution is empirically described as the ability to discriminate
between two points. If an object is just below the level of resolution, the peaks generated
by the two points will make the object appear to be a single point.
Abbe's equation. Resolution in a perfect optical system can be described

mathematically by Abbe's equation. In this equation:
d = _0.612 * 
n sin 
where:
d = resolution
 = wavelength of imaging radiation
n = index of refraction of medium between point source and lens, relative to free space
 = half the angle of the cone of light from specimen plane accepted by the objective
(half aperture angle in radians)
n sin  is often expressed as NA (numerical aperture)

This is the diffraction-limited resolution of an optical system. If all aberrations and

distortions are eliminated from the optical system, this will be the limit to resolution. If
aberrations and distortions are present, they will determine the practical limit to
resolution.

7. Elasticity
Modulus of Elasticity
Rate of change of strain as a function of stress. The slope of the straight-line portion of
a stress-strain diagram. Tangent modulus of elasticity is the slope of the stress-strain
diagram at any point. Secant modulus of elasticity is stress divided by strain at any
given value of stress or strain. It also is called stress-strain ratio.
Tangent and secant modulus of elasticity are equal, up to the proportional limit of a
material. Depending on the type of loading represented by the stress-strain diagram,
modulus of elasticity may be reported as: compressive modulus of elasticity (or
modulus of elasticity in compression); flexural modulus of elasticity (or modulus of
elasticity in flexure); shear modulus of elasticity (or modulus of elasticity in shear);
tensile modulus of elasticity (or modulus of elasticity in tension); or torsional modulus
of elasticity (or modulus of elasticity in torsion). Modulus of elasticity may be
determined by dynamic testing, where it can be derived from complex modulus.
Modulus used alone generally refers to tensile modulus of elasticity. Shear modulus is
almost always equal to torsional modulus and both are called modulus of rigidity.
Moduli of elasticity in tension and compression are approximately equal and are
known as Young's Modulus. Modulus of rigidity is related to Young's Modulus by the
equation:
where E is Young's Modulus (psi), G is modulus of rigidity (psi) and r is Poisson's

ratio. Modulus of elasticity also is called elastic modulus and coefficient of elasticity.

The ratio of stress (nominal) to corresponding strain below the proportional limit of a
material. It is expressed in force per unit area, usually pounds per square inch or
kilograms-force per square centimeter. The strain may be a change in length. Young's
modulus; a twist or shear, modulus of rigidity or modulus of torsion; or a change in
volume, bulk modulus. In the SI, all types of moduli of elasticity are reported in
pascals, the conversion factor being one psi = 6.894 757 E+3 pascals. Also known as
ELASTIC MODULUS and YOUNG'S MODULUS.’
Hooke’s Law:
The extension is proportional to the force or tension in a wire if the proportional limit
is not exceeded.
This is the fundamental law of elasticity and was first enunciated in the year 1679 by
Robert Hooke in the concise form ‘ Ut tensio sic vis ‘ i.e As the tension, so the strain. It
may, however be formally stated thus:
‘Provided the strain is small, the stress is proportional to the strain’ It follows,
therefore, that if the strain be small, the ratio between the stress and strain is a constant,
called modulus of elasticity or coefficient of elasticity (E). Thus E = stress/ strain or
stress per unit strain.
Since stress has the same dimensions and units as those of pressure and strain is a mere
number having no dimensions, the dimensions and the units of the modulus or
coefficient of elasticity are the same as those of stress and pressure
The extension of the wire is due to displacement of its molecules from their mean
positions. So the law shows that when a molecule of the metal is slightly displaced from
its mean position, the restoring force on the molecule is proportional to the displacement.
We may therefore conclude that the molecules of a solid metal are undergoing a simple
harmonic motion. Up to the elastic limit of the energy gained or stored by a stretched
wire is molecular potential energy, which is recovered when the load is removed.
The measurement also shows that it would be dangerous to load the wire with weights
greater than the magnitude of the elastic limit, because the wire then suffers a permanent
strain. Similar experiments in the research laboratory enable scientists to find the
maximum load, which a steel bridge, for example, should carry for safety. Rubber

samples are also subjected to similar experiments, to find the maximum safe tension in
rubber belts used in machinery.
Bending of Beams:
A rod or a bar of a circular or rectangular cross section, with its length very much greater
than its thickness is called a beam. Now a beam may just rest on a support, like a knife-
edge or have a small part of it firmly clamped or built into a wall at either end. In the
former case, it is called a supported beam or usually, simply a fixed beam.
In a supported beam, obviously the support can merely exert a force on the beam but in a
fixed beam, it can exert a couple on it.
Bending of a Beam:
Suppose we have a beam, of a rectangular cross section, say, fixed at one end and loaded
at the other so as to be bent a little, as shown in Fig.13.1, with its upper surface
becoming slightly convex and the lower one concave. All the longitudinal filaments in
the upper half of the beam thus gets extended or lengthened, and therefore under tension,
and all those in the lower half get compressed or shortened and therefore under pressure.
These extensions and compressions increase progressively as we proceed away from the
axis on either side ( as shown in a rough and ready manner by the lengths of the arrows),
so that they are the maximum in the uppermost and the lowermost layers of the beam
respectively.
Obviously there must be a layer between the uppermost and the lowermost layers where
the extensions in the upper half change sign to become compressions in the lower half.
In this layer or plane, which is perpendicular to the section of beam containing the axis,
the filaments never get extended nor compressed, i.e. retain their original lengths. This
layer is therefore called the Neutral surface of the beam, shown shaded in the Fig.13.2
If the material of the beam be homogenous and isotropic and the bending uniform, the
longitudinal filaments all get bent into circular arcs in planes parallel to the plane of
symmetry ABCD of the beam, which is also, therefore, the plane of bending. The centers
of curvature of all theses arcs lie on a straight line PQ, per-perpendicular to the plane of
bending (or the plane of symmetry), which is referred to as the axis of bending.
The line of intersection (OO’) of the plane of bending and the neutral surface (which are
clearly perpendicular to each other) is called the neutral axis of the beam.

The extension and compression of the filaments are directly proportional to their
respective distances from this axis.

8. Magnetism
The Earth a magnet
We probably know that the Earth is stratified; a section is pictured here. In radius it is
composed of layers having different chemical composition and different physical
properties. The crust of the Earth has some permanent magnetization, and the core of the
Earth, the outer part of which is liquid iron and the inner of which is solid iron, generates
its own magnetic field, sustaining the main part of the field we measure at the surface.
So we could say that the Earth is, therefore, a ‘magnet’. But there is no giant bar magnet
near the Earth’s center, despite the depictions you may have seen in elementary
textbooks on geology and geophysics. Permanent magnetization cannot occur at high
temperatures, like temperatures above 650 degrees centigrade or so, when the thermal
motion of atoms becomes sufficiently vigorous to destroy the ordered orientations
needed to establish permanent magnetization. The core of the Earth has a temperature of
several thousand degrees, and therefore it is not permanently magnetized.
Applying Knowledge: The Earth as a Magnet

 Draw Figure 2 on the board or overhead.
 Ask students to imagine that the Earth behaves like a magnet.
 Have students copy your version of figure 2 on a sheet of paper.
 Remind them of their results of working with the magnet
 ask them to draw the pattern that iron filings might make if they were sprinkled
all around the Earth
 Walk around the room and look at student work.
 Ask for volunteers to hold up their drawings and explain their ideas to the class.

Figure 2 Figure 3
Earth
 Next, draw Figure 3 on the board or overhead. Explain to students that this is the
actual shape of the magnetic field that surrounds Earth. Ask students to
hypothesize what might account for the distortion of the Earth’s magnetic field.
List their ideas on the board.
Review student suggestions and concluded that:

 The Earth has a magnetic field that is distorted by the solar wind
 The magnetic field is constantly changing in response to activity of the Sun.

Extend Student Inquiry

If students work quickly or a teacher requests ideas for follow-up activities, consider the
options below:
Have students repeat their observations using magnets of various shapes and strengths.
Have students collect their own iron filings by dragging a magnet through sand.
Students may do this with a box of sand or outside in places such as dunes, the seashore,
or a desert.
Identify, describe, and explain the existence, or lack of, a magnetic field around the other
planets in our solar system.
Magnetic Field Characteristics
Magnetic Field In and Around a Bar Magnet
As discussed previously a magnetic field is a change in energy within a volume of space.

The magnetic field surrounding a bar magnet can be seen in the magnetograph below. A
magnetograph can be created by placing a piece of paper over a magnet and sprinkling
the paper with iron filings. The particles align themselves with the lines of magnetic
force produced by the magnet. The magnetic lines of force show where the magnetic
field exits the material at one pole and reenters the material at another pole along the
length of the magnet. It should be noted that the magnetic lines of force exist in three-
dimensions but are only seen in two dimensions in the image.

It can be seen in the magnetograph that there are poles all along the length of the magnet
but that the poles are concentrated at the ends of the magnet. The area where the exit
poles are concentrated is called the magnet's north pole and the area where the entrance
poles are concentrated is called the magnet's south pole.
Magnetic Fields in and around Horseshoe and Ring Magnets
Magnets come in a variety of shapes and one of the more common is

the horseshoe (U) magnet. The horseshoe magnet has north and south
poles just like a bar magnet but the magnet is curved so the poles lie in
the same plane. The magnetic lines of force flow from pole to pole just
like in the bar magnet. However, since the poles are located closer
together and a more direct path exists for the lines of flux to travel
between the poles, the magnetic field is concentrated between the poles.
If a bar magnet was placed across the end of a horseshoe magnet or if a magnet was
formed in the shape of a ring, the lines of magnetic force would not even need to enter
the air. The value of such a magnet where the magnetic field is completely contained

with the material probably has limited use. However, it is important to understand that
the magnetic field can flow in loop within a material when the concept of circular
magnetism is later covered.
General Properties of Magnetic Lines of Force
Magnetic lines of force have a number of important properties, which include:
 They seek the path of least resistance between opposite magnetic poles.
In a single bar magnet as shown to the right, they attempt to form closed
loop from pole to pole.
 They never cross one another.
 They all have the same strength.
 Their density decreases (they spread out) when they move from an area
of higher permeability to an area of lower permeability.
 Their density decreases with increasing distance from the poles.

They are considered to have direction as if flowing, though no actual movement occurs.
They flow from the south pole to the north pole within the material and north pole to
south pole in air
Bar Magnet
The lines of magnetic field from a bar magnet form closed lines. By convention, the field
direction is taken to be outward from the North pole and in to the South pole of the
magnet. Permanent magnets can be made from ferromagnetic materials
All magnetic objects produce invisible lines of force that extend between the poles of the
object. An easy way to visualize this is to spread iron filings on a sheet of paper and
place a bar magnet under the paper. The iron filings will arrange themselves around the
magnet and along the magnetic field lines.

In the simplest terms, Earth can be thought of as a dipole (2-pole) magnet. Magnetic
field lines radiate between Earth's north and south magnetic poles just as they do
between the poles of a bar magnet. Charged particles become trapped on these field lines
(just as the iron filings are trapped), forming the magnetosphere.
All magnetic objects produce invisible lines of force that extend between the poles of the
object. An easy way to visualize this is to spread iron filings on a sheet of paper and
place a bar magnet under the paper. The iron filings will arrange themselves around the
magnet and along the magnetic field lines.
In the simplest terms, Earth can be thought of as a dipole (2-pole) magnet. Magnetic
field lines radiate between Earth's north and south magnetic poles just as they do
between the poles of a bar magnet. Charged particles become trapped on these field lines
(just as the iron filings are trapped), forming the magnetosphere.
Earth's magnetic field lines are not as symmetrical as those of the bar magnet. The
impact of the solar wind causes the lines facing sunward to compress, while the field
lines facing away from the Sun stream back to form Earth's magnetotail. The
magnetosphere extends into the vacuum of space from approximately 80 to 60,000
kilometers (50 to 37,280 miles) on the side toward the Sun, and trails out more than
300,000 kilometers (186,500 miles) away from the Sun.
The cause of Earth's magnetic field (the surface magnetic field) is not known for certain,
but is possibly explained by dynamo theory. The magnetic field extends several tens of
thousands of kilometers into space.
The field is approximately a magnetic dipole, with one pole near the geographic north
pole and the other near the geographic south pole. An imaginary line joining the
magnetic poles would be inclined by approximately 11.3° from the planet's axis of
rotation. The strength of the field at the Earth's surface at this time ranges from less than

30 microtesla (0.3 gauss) in an area including most of South America and South Africa
to over 60 microtesla (0.6 gauss) around the magnetic poles in northern Canada and
south of Australia, and in part of Siberia.
The field is similar to that of a bar magnet, but this similarity is superficial. The
magnetic field of a bar magnet, or any other type of permanent magnet, is created by the
coordinated motions of electrons (negatively charged particles) within iron atoms. The
Earth's core, however, is hotter than 1043 K, the temperature at which the orientations of
electron orbits within iron become randomized. Such randomization tends to cause the
substance to lose its magnetic field. Therefore the Earth's magnetic field is caused not by
magnetised iron deposits, but mostly by electric currents (known as telluric currents).
Another feature that distinguishes the Earth magnetically from a bar magnet is its
magnetosphere. At large distances from the planet, this dominates the surface magnetic
field. In addition, the magnetized elements within the planetary core are undergoing
rotation and are not static.
Intensity of magnetic field
The magnitude of the magnetic field at any point.

Strictly speaking, the magnetic field is the set of all values of the magnetic field
intensity, but magnetic field and magnetic field intensity (as well as magnetic field
strength and magnetic vector) are used more or less interchangeably. The trend is to use
magnetic field both for the field taken as a whole and for its value at any point, context
being sufficient to determine precise meaning.
Magnetic field reversals

The Earth's magnetic field reverses at intervals, ranging from tens of thousands to
hundreds of thousands of years. It is believed that this last occurred some 600,000 years
ago (Comins - DEU p.84). Past field reversals are recorded in the "frozen" magnetic
domains of solidified lava that has welled up along spreading ocean floor ridges; since

the sea floor spreads at a fairly constant rate, this results in broad "stripes" of sea floor
with opposite magnetic fields.
For the past 2000 years Earth's magnetic field has been weakening. At the going rate of
decay, the north-south magnetic dipole—generated within the convecting metallic fluid
of Earth's outer core—would totally vanish, perhaps reversing polarity in the next 2000
years.
This scenario of a coming attempt by Earth's magnetic field to reverse its polarity is
suggested by direct observation of the field since the 19th century and laboratory
investigation of historic lavas and other fired materials that record the ambient field
while cooling.
The ongoing weakening of the field does not ensure that a reversal will occur. After all,
Earth's dipole reverses direction only on occasion, currently at a rate of a few times each
million years. How a change in polarity is actually approached and, moreover, the degree
to which such a process can be predicted, are unclear. Nonetheless, a significant step
toward such an understanding may have been made through investigations of the
ancient, or paleomagnetic field recorded in Ocean Drilling Project (ODP) marine
sediment cores.
Such a process brings to mind an early statistical reversal model developed in the late
1960s. It proposed that field reversal could only occur when the dipole field can interact
with the typically weaker, more complex nondipole part of the field. According to the
model, this occurs when dipole strength, assumed to vary in a sinusoidal wavelike
manner, is at an intensity low. The newly reported sediment data, however, suggest that
any substantial strengthening of the dipole requires an attempt by the field to reverse
and, hence, links field generation to the reversal process.

Earth is currently in the Brunhes normal polarity chron. By presenting globally

distributed paleomagnetic intensity data from the Pacific, Atlantic, and Indian Oceans
crossing the boundary of the last known reversal (the 780,000-year-old Matuyama-
Brunhes reverse-to-normal polarity change), the investigators in Paris successfully tested
their idea. These comparisons add credence to the claim that the asymmetric "saw-
toothed" decay-and-recovery intensity pattern originally observed is a worldwide
phenomenon linked to the dipole field. Partially compatible with this general picture are
paleomagnetic data obtained from sequences of lava flows that recorded field behavior
during reversals. These studies show that just following a polarity change the dipole field
can strengthen to values well beyond that experienced prior to the flip. However, no
long-term gradual field decay has been noted from lavas.
If the ocean sediment data are correct, then the process of polarity reversal of the
geomagnetic field cannot be considered purely random, even though the time spanned
between reversals itself appears to be random. Indeed, the researchers in Paris contend
that the sediment data indicate a link between the time needed by the dipole field to
weaken to the point of possible reversal and the degree of revitalization of the field
following the previous reversal.
On another side of the issue is a statistical model developed in 1993. It is based on the
known timing of field reversals over the past 100 million years and suggests that
following a reversal there exists a 5,000 year-long dead-time during which there is zero
probability of another occurrence. Then, according to the model, the probability of
reversal steadily increases with time, but only for some 45,000 years. Although rooted in
observed long-term field behavior, this model is not fully compatible with the French
researchers' claim that the duration of a given polarity interval may be virtually
determined by the initial field strength following reversal.
The sediment-derived finding of such a remarkable saw-toothed paleointensity pattern

apparently linked to the reversal process has stimulated much discussion and debate
among experimentalists and theorists who are exploring the source of Earth's magnetic

field. Alternative nongeomagnetic explanations based on the process by which sediments

become magnetized are also being sought. Whatever the outcome, this growing
controversy promises to significantly affect our perceptions regarding field reversal.
Whether we may one day understand the process sufficiently well to formulate
predictions about the next polarity change is yet to be determined.
Angle of Dip of earths Field
A dip circle is used to measure the angle between the direction of the Earth's magnetic
field and the horizontal at any point on the Earth's surface. Angle at a particular point on
the Earth's surface between the direction of the Earth's magnetic field and the horizontal.
It is measured using a dip circle, which has a magnetized needle suspended so that it can
turn freely in the vertical plane of the magnetic field. In the northern hemisphere the
needle dips below the horizontal, pointing along the line of the magnetic field towards its
north pole. At the magnetic north and south poles, the needle dips vertically and the
angle of dip is 90°. See also angle of declination.
Dip errors
Only near the equator is earth's magnetic field parallel to the earth's surface. In the
Northern hemisphere the field lines have a downward component and in the Southern,
upward. The angle that the magnetic field makes with the horizontal is called the dip
angle, and like variation, it varies from place to place and slowly with time. Charts of dip
and variation can be found on the net and or can be computed using published formulae.

If the needle of the compass were mounted so it could pivot freely about its center of
gravity in three dimensions, it would align with the magnetic field, pointing down (or
up) at the dip angle in the direction of local Magnetic North. Since the dip angle is not of
navigational interest, the compass is constructed so that is constrained to rotate
essentially only in the horizontal plane.
In an aviation compass, this is done by lowering the CG below the pivot point and
making the assembly heavy enough that the vertical component of the magnetic force is
too weak to tilt it significantly out of the horizontal plane. The compass can then work
effectively at all latitudes without specific compensation for dip. However, close to the
magnetic poles, the horizontal component of the earth's field is too small to align the
compass, and the compass becomes of little use for navigation.
Because of this constraint, the compass only indicates correctly if its card is horizontal.
If it is tilted out of the horizontal, it will be affected by the vertical component of the
earth's field. This gives rise to two distinct errors.
dip, magnetic
A dip circle is used to measure the angle between the direction of the Earth's magnetic
field and the horizontal at any point on the Earth's surface.
(Image © Research Machines plc)
Angle at a particular point on the Earth's surface between the direction of the Earth's
magnetic field and the horizontal. It is measured using a dip circle, which has a
magnetized needle suspended so that it can turn freely in the vertical plane of the
magnetic field. In the northern hemisphere the needle dips below the horizontal, pointing
along the line of the magnetic field towards its north pole. At the magnetic north and
south poles, the needle dips vertically and the angle of dip is 90°. See also angle of
declination.

Domain Theory
A remarkable property of ferrimagnetic materials is not so much that they have a
spontaneous magnetization, but rather that their magnetization can be influenced by the
application of very low magnetic fields. Even the earth's field (50 T) can cause
magnetization changes even though the interatomic exchange forces responsible for the
spontaneous magnetization are equivalent to a field of about 1000 T, almost 100 million
times greater than the earth's field.
What allows this to occur is the fact that the sample is actually composed of small
regions called magnetic domains, within each of which the local magnetization is
saturated but not necessarily parallel. Domains are small (1-100's microns), but much
larger than atomic distances.
The existence of domains is hinted at by the observation that some magnetic properties,
and in particular, coercivity and remanence vary greatly with grain size. This is best
illustrated in the figure below, which shows the variation of Hc with grain size.

The magnetic behavior can be subdivided on the basis of grain size into four ranges
SPM: superparamagnetic
SD: single domain
PSD: pseudo-single domain
MD: multidomain
The maximum coactivity for a given material occurs within its SD range. For larger
grain sizes, coercivity decreases as the grain subdivides into domains. For smaller grain
sizes, coercivity again decreases, but this time due to the randomizing effects of thermal
energy.
Domains constitute a fundamental concept in magnetism. A ferro- or ferrimagnetic

material may be generally defined as one that possesses a spontaneous magnetization,
Ms, dependent on temperature, but only slightly dependent on applied field. The theory
of ferromagnetism, based on electronic exchange forces, predicts the magnitude of Ms,
but says nothing about the direction of Ms. Experimentally, it is observed that for a
homogeneous specimen at constant temperature, the magnitude of Ms is uniform but the
direction of Ms is in general not uniform from one region to another (on a scale of
microns to millimeters). Uniformity of direction is attained only by applying a large
enough field to drive the domains out of the sample, or by reducing the particle's
dimensions to small enough size to prevent domain formation.

Domains are formed for the following reason. Consider a large single crystal.
Suppose it is uniformly magnetized, and hence a single domain. Surface charges will
form on the ends due to the magnetization and are themselves a second source of a
magnetic field (the demagnetizing field). The energy associated with the surface charge
distribution is called the magnetostatic energy. It is just the volume integral of the field
over all space.
The magnetostatic energy can be approximately halved if the magnetization splits into
two domains magnetized in opposite directions. This brings (+) and (-) charges closer
together, thus decreasing the spatial extent of the demagnetizing field.
This subdivision into more and more domains can not continue indefinitely because the
transition region between domains (called a domain wall) requires energy to be produced
and maintained. Eventually an equilibrium number of domains will be reached for a
given particle size.

Domain walls are interfaces between regions in which the magnetization has different
directions. Within the wall, the magnetization must change direction from that in one
domain to that in the other domain. Domain walls have a finite width that is determined
principally by exchange and magnetocrystalline energy.
Let's consider a domain wall in which the magnetization changes by 180°. The change in
magnetization within the wall can be gradual as in (a) or abrupt as in (b).
The exchange energy acts to keep spins parallel and can be kept small if the
180°rotation takes place gradually, over many atomic units. Therefore, the exchange
energy is small in (a) but large in (b).
However, the spins within the wall are no longer aligned along an easy axis of
magnetization. This produces an anisotropy energy, which is high in (a) but low in (b).
The exchange energy tends to make the wall as wide as possible whereas the anisotropy
tends to make the wall as thin as possible. As a result of this competition between
exchange and anisotropy energies,the domain wall has a finite width (on the order of 100
nm) and surface energy.

The interplay between long range and short range effects results in the domain states
being grain-size dependent. In addition, the number of domains for a given grain size
depends on the magnitudes of the exchange, magnetocrystalline, and saturation
magnetization. As mentioned before, these constants are dependent on temperature as
well as composition. Hence domain states in different magnetic minerals (magnetite and
hematite) will have a different grain size dependence. The domain states will also vary
with temperature for a single grain size. However, as a rule of thumb, the larger the grain
size the more domains it contains.
Single Domain (SD)
As the grain size decreases, a critical size will be reached where the grain can no longer
accommodate a wall. Below this critical size, the grain contains a single domain (SD).
An SD grain is uniformly magnetized to its saturation magnetization.
SD grains are very important. To change the magnetization of a MD grain, all you need
to do is translate the domain wall, a energetically easy process, which can be
accomplished in relatively low fields. Thus MD grains are magnetically soft with low
values of coercivities and remanence
The Domain Theory of Magnetism

William Gilbert's report on his experiments with natural magnets, published in 1600,
probably represents the first scientific study of magnetism. In the years that followed,
discoveries by Coulomb, Oersted, and Ampere added to our knowledge of the behavior
of magnets and the nature of magnetic forces.
Physicists believe, however, that it is only within this century that they have begun to
understand the true nature of magnetism. The present view is that the magnetic
properties of matter are electric in origin and result from the movements of electrons
within the atoms of substances.

Since the electron is an electrically charged particle, this theory suggests that magnetism
is a property of a charge in motion. If so, we can account for the energy associated with
magnetic forces by using known laws of physics.
Two kinds of electron motion are important in this modern concept of magnetism. First,
an electron revolving about the nucleus of an atom imparts a magnetic property to the
atom structure. See Figure 19-2.
When the atoms of a substance are subjected to the magnetic force of a strong magnet,
the force affects this magnetic property, opposing the motion of the electrons. The atoms
are thus repelled by the magnet. This is diamagnetism. If the electron's only motion were
its movement about the nucleus, all substances would be diamagnetic.
Diamagnetic repulsion is quite feeble in its action on the total mass of a substance. The
second kind of motion is that of the electron spinning on its own, axis. Each spinning
electron acts as a tiny permanent magnet. Opposite spins are designated as + and - spins;
electrons spinning in opposite directions tend to form pairs and so neutralize their
magnetic character. See Figure 19-3 (above). The magnetic character of an atom as a
whole may be weak because of the mutual interaction between the electron spins.

Magnetic properties are associated with both kinds of electron motion. The atoms of
some substances may possess permanent magnet characteristics because of an
imbalance between orbits and spins. These atoms act like tiny magnets, called dipoles,
and are attracted by strong magnets.
Substances in which this attractive effect exceeds the diamagnetism common to all atoms
show the property of paramagnetism.
In the atoms of ferromagnetic substances there are unpaired electrons whose spins are
oriented in the same way. The common metals iron, cobalt, and nickel and the rare earth
elements gadolinium and dysprosium show strong ferromagnetic properties.
Some alloys of these and other elements, as well as certain metallic oxides called
ferrites, also exhibit strong ferromagnetic properties. The inner quantum levels, or shells,
of the atom structures of most elements contain only paired electrons. The highest
quantum level, or outer shell, of each of the noble gases (except helium) consists of a
stable octet of electrons made up of four electron pairs.
The atoms of other elements achieve this stable configuration by forming chemical
bonds. Only in certain transition elements that have incomplete inner shells do unpaired
electrons result in ferromagnetic properties. The electron configuration of the iron atom,
Figure 19-4, shows four unpaired electrons in the third principal quantum level. The
similarly oriented spins of these electrons, enhanced by the influence of nearby atoms in
the metallic crystal, account for iron's strong ferromagnetism.

From the preceding discussion, it would seem that every piece of iron should behave as a
magnet. However, such is not the case. Atoms are grouped in microscopic magnetic
regions called domains. The atoms in each domain are magnetically polarized parallel to
a crystal axis. In a polycrystalline specimen, ordinarily these axes (and the domains) are
oriented in all possible directions. The domains effectively cancel one another and the
net magnetism is essentially zero. In Figure 19-5 the polarity of each domain in an
unmagnetized material is represented by an arrow.
When a ferromagnetic material is placed in an external magnetic field, two effects occur.
The domains more favorably oriented in this magnetic field may increase in size at the
expense of less favorably oriented adjacent domains. Other domains may rotate in order
to become more favorably oriented with respect to the external field. The material
becomes magnetized.
If the domain boundaries remain extended to some degree even after the external
magnetizing force is removed, the material is said to be "permanently" magnetized.
When the direction of magnetization of a magnetic domain is rotated by an external
magnetic field, it must be understood that the material of the domain does not change its
position in the specimen. It is only its direction of magnetization that changes.

When the temperature of a ferromagnetic material is raised above a certain critical

value, the domain regions disappear and the material becomes paramagnetic. This
temperature is known as the Curie point. It is usually lower than the melting point of the
substance. The Curie points for some ferromagnetic substances are given in Table 19-1.
When a single crystal of iron is sprinkled with colloidal particles of iron oxide, the
microscopic domains become visible. Using this technique, physicists are able to
photograph magnetic domains and observe the effects of external magnetic fields on
them.
Diamagnetism
When a diamagnetic material is placed near a magnet, it will be repelled from the region
of greater magnetic field, just opposite to a ferromagnetic material. It is exhibited by all
common materials, but is very weak. People and frogs are diamagnetic. An interesting
experiment showing this is at http://www.sci.kun.nl/hfml/froglev.html where a frog is
levitated at the top of a very strong electromagnet. Metals such as bismuth, copper, gold,
silver and lead, as well as many nonmetals such as graphite, water and most organic
compounds are diamagnetic.

For example, here are some photos of a very strong neodymium-iron-boron magnet
sitting in a dish with a shallow amount of water covering the magnet. Looking at the
reflection of the light above the sink off of the surface of the water, you can see how the
reflection is distorted because the water is concave just above the magnet, and flat
everywhere else. (The fuzzy object in the two photos on the right is the magnet; the
camera is focused on the reflected light.) This is because the magnet has pushed the
water away since water is repelled by strong magnetic fields.
Here's a fixture using diamagnetic disks that you can purchase and experiment with.
Diamagnetic
Diamagnetic is a classification used to describe materials that line up at right angles to a

nonuniform magnetic field and are slightly repelled by that field. Diamagnetism occurs
as a result of a magnetic field's interference with the motion of electrons orbiting the
atoms or molecules of an element or compound. When matter is placed in a magnetic
field, the magnetic force acts upon the moving electrons in the matter, causing the
electrons to be deflected in ways that speed up some electrons and slow down others.
This movement of the electrons interferes with the motion of the magnetic field, so the
atoms internally oppose the field. This causes the material to be slightly repelled by the
magnetic field.
Diamagnetism is characteristic of elements and compounds that possess complete sets of

valence electrons--meaning that all their electrons are paired. Electrons orbit an atom
while spinning on their own axis. If a spinning electron is orbiting singly, this charge
movement generates a magnetic field. However, two paired electrons have opposite
spins--meaning that the magnetic field generated by each electron is canceled out by the
magnetic field of the other electron. So when all of an atom's electrons are paired, the

opposite spins of the electrons cancel each other out and cause the element to have no
magnetic field. When this element is then placed in a magnetic field, it is repelled.
All materials demonstrate a degree of diamagnetism. In some substances, however, other

magnetic properties may mask its diamagnetic properties. Paramagnets, another
classification of matter, are one example of one magnetic property masking the
diamagnetic property. Paramagnetism generally occurs in elements or compounds
possessing unpaired electrons. Many compounds containing iron, palladium, platinum,
and the rare earth elements have single electrons that generate a small magnetic field. In
these cases, the atom acts as a small permanent magnet. If a substance containing such
atoms is placed in a magnetic field, the field of the atom aligns with the field of the
applied magnetic field and causes the atom to be slightly attracted to that magnetic field.
This attraction to an applied magnetic field is known as paramagnetism
Paramagnetism
When a paramagnetic material is placed near a magnet, it will be attracted to the region
of greater magnetic field, like a ferromagnetic material. The difference is that the
attraction is weak. It is exhibited by materials containing transition elements, rare earth
elements and actinide elements. Liquid oxygen and aluminum are examples of
paramagnetic materials.
Ferromagnetism
When a ferromagnetic material is placed near a magnet, it will be attracted toward the
region of greater magnetic field. This is what we are most familiar with when our
magnet picks up a bunch of paperclips. Iron, cobalt, nickel, gadolinium, dysprosium and
alloys containing these elements exhibit ferromagnetism because of the way the electron
spins within one atom interact with those of nearby atoms. They will align themselves,
creating magnetic domains forming a temporary magnet. If a piece of iron is placed
within a strong magnetic field, the domains in line with the field will grow in size as the
domains perpendicular to the field will shrink in size.

Form of magnetism that can be acquired in an external magnetic field and usually
retained in its absence, so that ferromagnetic materials are used to make permanent
magnets. A ferromagnetic material may therefore be said to have a high magnetic
permeability and susceptibility (which depends upon temperature). Examples are iron,
cobalt, nickel, and their alloys.
Ultimately, ferromagnetism is caused by spinning electrons in the atoms of the material,

which act as tiny weak magnets. They align parallel to each other within small regions of
the material to form domains, or areas of stronger magnetism. In an unmagnetized
material, the domains are aligned at random so there is no overall magnetic effect. If a
magnetic field is applied to that material, the domains align to point in the same
direction, producing a strong overall magnetic effect. Permanent magnetism arises if the
domains remain aligned after the external field is removed. Ferromagnetic materials
exhibit hysteresis
Ferromagnetic materials are the most magnetically active substances in the world, and
so they have very high magnetic susceptibilities, ranging from 1000 up to 100,000.
These materials are made of atoms with permanent dipole moments, and when these
materials form solids by exchanging electrons to make chemical bonds, something
special happens. If the atoms are of the right type and if the bond lengths are right, the
electrons discover that they can place the system in a state of lower energy by having
neighboring atomic dipole moments aligned with each other. (This sentence probably
makes no sense to you, but it is the best we can do. That's just the way quantum
mechanics is.) If the entire sample were to be made of aligned dipoles, however, a strong
magnetic field would be created, and this would be a state of high energy. So the system
compromises. It makes microscopic regions in which billions of dipoles are aligned,
satisfying the demands of most of the electron bonds. But the alignment directions of the
separate regions are random throughout the sample, making a very weak net magnetic
field. These regions are called magnetic domains, and their behavior gives
ferromagnetic materials their distinctive properties.

For example, if a ferromagnetic material is heated to too high a temperature, it ceases to

be ferromagnetic. The reason is that above a certain critical temperature, called the Curie
temperature, the thermal motion of the atoms is so violent that the electrons in the bonds
are no longer able to keep the dipole moments aligned. When this happens, the
ferromagnetic material changes into a paramagnetic material with the usual weak
magnetism.
The domains also make permanent magnets possible. If a ferromagnetic sample is placed
in a strong magnetic field, the domains can be forced to coalesce into large domains
aligned with the external field. When the external field is removed, the electrons in the
bonds maintain the alignment and the magnetism remains. This means that
ferromagnetic materials can remember their past magnetic history. This property of
magnetic memory is called hysteresis and lies at the heart of audio tape, video tape,
and magnetic disk storage for computers. The recording head of a tape recorder, or the
write head of a disk drive, applies a field that magnetizes a small portion of the tape or
disk. The magnetism in each portion remains until another magnetic field changes it.
When each magnetized section is moved under the playback head of a tape player, or the
read head of a disk drive, the moving magnetic field induces small currents which are
amplified and turned into either music or data bits. If the domains were unable to
remember the field that had been applied to them, none of this would be possible
Magnetically Hard and Soft Iron:
The opposite extreme to a ``hard'' ferromagnetic material, which can maintain a large
remnant magnetization in the absence of external fields, is a ``soft'' ferromagnetic
material, for which the remnant magnetization is relatively small. Let us consider a
somewhat idealized situation in which the remnant magnetization is negligible. In this
situation there is no hysteresis, so the - relation for the material reduces to
(591)

where is a single valued function. The most commonly occurring ``soft''
ferromagnetic material is soft iron (i.e., annealed, low impurity iron).
Consider a sphere of soft iron placed in an initially uniform external field .
The and fields inside the sphere are most easily obtained by taking the solutions
(3.166) and (3.167) (which are still valid in this case), and superimposing on them the
uniform field . We are justified in doing this because the equations which govern
magnetostatic problems are linear. Thus, inside the sphere we have
(592)
(593)
Combining Eqs. (3.169) and (3.170) yields
(594)
with
(595)

where, in general, . Clearly, for a highly permeable material (i.e., ,
which is certainly the case for soft iron) the magnetic field strength inside the sphere is
approximately three times that of the externally applied field. In other words, the
magnetic field is amplified inside the sphere.
The amplification of the magnetic field by a factor three in the high permeability limit is
specific to a sphere. It can be shown that for elongated objects (e.g., rods), aligned along
the direction of the external field, the amplification factor can be considerably larger
than this.
It is important to realize that the magnetization inside a ferromagnetic material cannot
increase without limit. The maximum possible value of is called the saturation
magnetization, and is usually denoted . Most ferromagnetic materials saturate when
they are placed in external magnetic fields whose strengths are greater than, or of order,
one tesla. Suppose that our soft iron sphere first attains the saturation magnetization
when the unperturbed external magnetic field strength is . It follows from (3.170b)
and (3.171) (with ) that
(596)
inside the sphere, for . In this case, the field amplification factor is

(597)
Thus, for the amplification factor approaches unity. We conclude that if a
ferromagnetic material is placed in an external field which greatly exceeds that required
to cause saturation then the material effectively loses its magnetic properties, so that
. Clearly, it is very important to avoid saturating the soft magnets used to channel
magnetic flux around transformer circuits. This sets an upper limit on the magnetic field
strengths which can occur in such circuits
Magnetic flux
is a measure of quantity of magnetism, taking account of the strength and the extent of a
magnetic field. The flux through an element of area perpendicular to the direction of
magnetic field is given by the product of the magnetic field density and the area element.
More generally, magnetic flux is defined by a scalar product of the magnetic field
density and the area element vector. Gauss's law for magnetism, which is one of the four
Maxwell's equations, states that the total magnetic flux through a closed surface is zero.
This law is a consequence of the empirical observation that magnetic monopoles do not
exist or are not measureable. The SI unit of magnetic flux is the weber, and the unit of
magnetic flux density is the weber per square meter.
In symbols, this means:

where is the magnetic flux and B is the magnetic flux density.
We know from Gauss's law that
This equation, in combination with the divergence theorem, provides the following
result:
In other words, the magnetic flux through any closed surface must be zero.
By way of contrast, Gauss's law for electric fields, another of Maxwell's equations, is
indicating the presence of electric monopoles. By the same token, there are no magnetic
monopoles.
The direction or vector of the magnetic flux is by definition from the south to the north
pole of a magnet (within the magnet). Outside of the magnet, the field lines will go from
north to south.
A change of magnetic flux in a spool of electrical conductive wire called a solenoid will
cause an electric current in the spool. This is the basis of the production of electricity.
When turned around, that is, running a current through a spool, a magnetic flux will be

produced in the spool. This is electromagnetism. The relationship is given by Faraday's

Law:
Very qualitatively, flux is a measure of how much of a vector field passes

perpendicularly through a given area. A vector field is simply a vector quantity that has
possibly different values (in magnitude and/or direction) at different points in space.
Examples of vector fields are the velocity of air molecules in a room or of water
molecules in a stream, the acceleration due to gravity at various locations on and above
the earth’s surface, the electric field surrounding a charged balloon, or the magnetic field
in the region around a solenoid with current flowing through its windings. (This
explains why the latter two examples are called fields; we can also speak of the
gravitational field around the earth, which is equivalent to the third example listed.) In
all of these cases, the vector quantity being described varies with position in space.
In this lecture we will be talking about the flux of the magnetic field. However, it will
be important to keep in mind that we can speak just as well of the flux of air molecule
velocities or the flux of the electric field as we can of the flux of the magnetic field. (See
the homework at the end of this lecture for specific examples.)
Again, magnetic flux is a measure of how much magnetic field passes perpendicularly
through a given area. Consider the flat loop of area A formed by N turns of wire shown
in Fig. 18.2.

Figure 18.2
A uniform magnetic field B passes perpendicularly through the turns of wire forming the
loop. The net flux of the magnetic field B, denoted B ( is the capital Greek letter
phi), through the area A in this case is given by
We see that, if we increase the area A or the magnetic field magnitude B or the number
of turns N, we will increase the net magnetic flux B. All of these should make perfect
sense in light of our qualitative understanding of flux. (The net flux is the total flux
through all of the loops of wire, if there are more than one.)
Consider now the situation shown in Fig. 18.3.
Figure 18.3

Again, we have N turns of wire defining the area A, and again the magnetic field is
uniform. But this time the magnetic field does not point perpendicular to the area A, but
rather makes an angle  with the perpendicular direction to the loop (the normal
direction to the plane of the loop – do you remember the normal force and what direction
it pointed?) Figure 18.4 shows the same set-up as in Fig. 18.3, but now shows just one
magnetic field vector (so that we can concentrate on just it and not worry about all of the
others cluttering things up). Any vector can be broken down into two perpendicular
components (like x and y components). We show in Fig. 18.4 the magnetic field vector
broken down into components perpendicular and parallel to the plane of the loop.
Figure 18.4
From our qualitative understanding of flux, we know that only the perpendicular part of
the field will contribute to B. (The parallel component just runs along the surface, and
does not pass through it at all.) From a straight-forward extension of Eq. (18.1), we get
that, in this case, the net magnetic flux through the loop is given by
This is the mathematical definition of magnetic flux that we will be using in this course.
(This definition assumes that the component of magnetic field passing perpendicular to
the area is constant over the whole area, and that the area is flat. If this is not the case,
then calculus would have to be used to find the flux. We won’t worry about that here!)
Note that in the set-up of Fig. 18.4, Eq. (18.2) becomes

For any given set-up, you must first find the component of the magnetic field
perpendicular to the area under consideration, and then multiply that component by the
area and the number of turns in order to find the flux through the loop.
The units of magnetic flux must be the units of number of turns (no units), times the
units of magnetic field (T) times the units of area (m2).
Units of magnetic flux: T m2 .
flux, magnetic, in physics, term used to describe the total amount of magnetic field in a
given region. The term flux was chosen because the power of a magnet seems to “flow”
out of the magnet at one pole and return at the other pole in a circulating pattern, as
suggested by the patterns formed by iron filings sprinkled on a paper placed over a
magnet or a conductor carrying an electric current. These patterns are called lines of
induction. Although there is no actual physical flow, the lines of induction suggest the
correct mathematical description of magnetism in terms of a field of force. The lines of
induction originate on the north pole of the magnet and end on the south pole; their
direction at any point is the direction of the magnetic field, and their density (the number
of lines passing through a unit area) gives the strength of the field. Near the poles where
the lines converge, the field and the force it produces are large; away from the poles
where the lines diverge, the field and force are progressively weaker.
Magnetic Flux through a Surface
 Magnetic flux measures the number of magnetic field lines that pass
perpendicularly through a surface.

SI: Weber = Wb =T.m 2
Gauss's Law for Magnetism
 The magnetic flux into a surface is equal to the magnetic flux out of a surface.
 This law is equivalent to the experimental observation that there are no magnetic
monopoles from which the magnetic field lines can originate.
 Compare this with Gauss' Law for the Electric Flux.
Flux Through a Loop

 When both the direction and magnitude of the magnetic field does not change
across the surface bounded by a closed loop, the magnetic flux can be expressed
as:
 = Angle between direction of B and A.
Direction of :

 The direction of is perpendicular to the surface of the loop. is normal to the

surface formed by the closed loop.
Thus the contribution to magnetic flux for a given area is equal to the area times the
component of magnetic field perpendicular to the area. For a closed surface, the sum of
magnetic flux is always equal to zero (Gauss' law for magnetism). No matter how small
the volume, the magnetic sources are always dipole sources (like miniature bar magnets),
so that there are as many magnetic field lines coming in (to the south pole) as out (from
the north pole).
Magneto motive force
Magnetomotive force, also known as magnetic potential, is the property of certain

substances or phenomena that gives rise to magnetic fields. Magnetomotive force is
analogous to electromotive force or voltage in electricity.
The standard unit of magnetomotive force is the ampere-turn (AT), represented by a

steady, direct electrical current of one ampere (1 A) flowing in a single-turn loop of
electrically conducting material in a vacuum. Sometimes a unit called the gilbert (G) is
used to quantify magnetomotive force. The gilbert is defined differently, and is a slightly
smaller unit than the ampere-turn. To convert from ampere-turns to gilberts, multiply by
1.25664. Conversely, multiply by 0.795773.
Although the standard definition of magnetomotive force involves current passing
through an electrical conductor, permanent magnets also exhibit magnetomotive force.
The same is true for planets with magnetic fields, such as the Earth, Jupiter, Saturn,
Uranus, and Neptune. The Sun also generates magnetomotive forces, particularly in the
vicinity of sunspots

A magnetic circuit consisting of a coil wound on either a magnetic or non-magnetic

former can be compared with the electric circuit. In the electric circuit:
Current = emf / resistance
or
I = E/R
In the magnetic circuit:
Flux = mmf / Reluctance
or
(phi) = F/Rm
Magnetomotive force is measured in amperes, A, and is produced by the current in the

magnetizing current where:
mmf = NI
where,
mmf is the magnetomotive force in Amperes, A.
N is the number of turns.
I is the magnetizing current in Amperes, A.
Permeability

This information is most easily obtained from the magnetization curve. Figure MPC
shows the permeability (in black) derived from the magnetization curve (in colour) using
equation MPD. Note carefully that permeability so defined is not the same as the slope
of a tangent to the B-H curve except at the peak (around 80 A m -1 in this case). The latter
is called differential permeability, μ′ = dB/dH.
If you use a core with a high value of permeability then fewer turns will be required to
produce a coil with a given value of inductance. You can understand why by
remembering that inductance is the ratio of flux to current. For a given core B is
proportional to flux and H is proportional to the current so that inductance is also
proportional to μ: the ratio of B to H.
Unlike electrical conductivity, permeability is often a highly non-linear quantity. Most

coil design formulæ, however, pretend that μ is a linear quantity. If you were working at
a peak value of H of 100 A m-1, for example, then you might take an average value for μ
of about 0.006 H m-1. This is all very approximate, but you must accept inaccuracy if you
insist on treating a non-linear quantity as though it was actually linear.
This form of permeability, where μ is written without a subscript, is known in SI

parlance as absolute permeability. It is seldom quoted in engineering texts. Instead a
variant is used called relative permeability described next.
Relative permeability

Relative permeability
Quantity name Relative permeability
Quantity symbol μr
Unit symbols dimensionless
Relative permeability is a very frequently used parameter. It is a variation upon

'straight' or absolute permeability, μ, but is more useful to you because it makes clearer
how the presence of a particular material affects the relationship between flux density
and field strength. The term 'relative' arises because this permeability is defined in
relation to the permeability of a vacuum, μ0
Equation
μr = μ / μ0
MPE
For example, if you use a material for which μr = 3 then you know that the flux density
will be three times as great as it would be if we just applied the same field strength to a
vacuum. This is simply a more user friendly way of saying that μ = 3.77×10 -6 H m-1.
Note that because μr is a dimensionless ratio that there are no units associated with it.
Many authors simply say "permeability" and leave you to infer that they mean relative
permeability. In the CGS system of units these are one and the same thing really. If a
figure greater than 1.0 is quoted then you can be almost certain it is μr.
Approximate maximum permeabilities

Material μ/(H m-1) μr Application
Ferrite U 60 1.00E-05 8 UHF chokes
Ferrite M33 9.42E-04 750 Resonant circuit RM cores
Nickel (99%
7.54E-04 600 -
pure)
Ferrite N41 3.77E-03 3000 Power circuits
Iron (99.8% pure) 6.28E-03 5000 -
Ferrite T38 1.26E-02 10000 Broadband transformers
Silicon GO steel 5.03E-02 40000 Dynamos, mains transformers
supermalloy 1.26 1000000 Recording heads

Note that, unlike μ0, μr is not constant and changes with flux density. Also, if the
temperature is increased from, say, 20 to 80 centigrade then a typical ferrite can suffer a
25% drop in permeability. This is a big problem in high-Q tuned circuits.
Another factor, with steel cores especially, is the microstructure, in particular grain
orientation. Silicon steel sheet is often made by cold rolling to orient the grains along the
laminations (rather than allowing them to lie randomly) giving increased μ. We call such
material anisotropic.
Before you pull any value of μ from a data sheet ask yourself if it is appropriate for your
material under the actual conditions under which you use it. Finally, if you do not know
the permeability of your core then build a simple circuit to measure it.
Effective permeability
Effective permeability
Quantity name Effective permeability
Quantity symbol μe
Unit symbols dimensionless *
Effective permeability is seen in some data sheets for cores which have air gaps. Coil
calculations are easier because you can simply ignore the gap by pretending that you are
using a material whose permeability is lower than the material you actually have.
* Effective permeability is usually relative to μ0.
Magnetic Field Strength

A magnetic field is a force field generated by moving electrical charges. An electric
current running through a loop of wire generates a magnetic field. The strength of the
field depends on the current and the area of the wire loop.
In the UK and the USA, the strength of the magnetic field is normally expressed as
Magnetic Field Strength (SI unit: amperes per metre, A/m), or Magnetic Flux Density
(B).

Consider the interior of a long wire coil in air.
Suppose that the linear current density in this coil is 1 ampere per metre along the coil
axis.Then the magnetic field strength in the interior of the coil is 1 A/m.
Magnetic field strength, also called magnetic intensity, is an expression of the force that
a magnetic field exerts on a theoretical unit magnetic pole in free space. This quantity
can be defined in various ways.
The unit of magnetic field strength in the International System of Units (SI) is the
ampere per meter (A/m). The oersted (Oe), which is the centimeter-gram-second (cgs)
unit, is sometimes used as an alternative unit in older documentation. To convert from
amperes per meter to oersteds, multiply by 0.012566. Conversely, multiply by 79.578.

Magnetic flux density and magnetic force are related to magnetic field strength. In
general, the magnetic field strength diminishes with increasing distance from the axis of
a magnetic dipole in which the flux field is stable. The function defining the rate at
which this field-strength decrease occurs depends on the geometry of the magnetic lines
of flux (the shape of the flux field).
A magnet produces a vector field, the magnetic field, at all points in the space around it.
It can be defined by measuring the force the field exerts on a moving charged particle,
such as an electron. The force (F) is equal to the charge (q) times the speed of the
particle times the magnitude of the field (B), or F = q*v x B, where the direction of F is
at right angles to both v and B as a result of the cross product. This defines the magnetic
field's strength and direction at any point.
A magnetic field can be created with moving charges, such as a current-carrying wire. A
magnetic field can also be created by the spin magnetic dipole moment, and by the
orbital magnetic dipole moment of an electron within an atom.
This is the Right Hand Rule for magnetic field from flowing current, and for magnetic
field in a coil.
When current flows in a wire, a magnetic field is created around the wire. To visualize
this, take your right hand, curl the fingers, and stick the thumb straight out. Now, point
your thumb in the direction of the current flowing in the wire (using conventional
current where current flows from the + end of a battery to the - end of the battery.
NOTE: electrons flow from the - end of a battery to the + end, and is called electron
current instead of conventional current). The direction your fingers are curved around

the wire is the direction of the magnetic field around the wire. For example, if the
current were coming straight out of this page toward you, your thumb would be pointing
toward you and your fingers would indicate a counter-clockwise direction to the
magnetic field around the wire.
If you have a coil of wire, simply curve the fingers of your right hand around the coil in
the same direction as the current is flowing. Your thumb will point out of the north
magnetic pole which the coil of wire will create.
By convention, we state that the magnetic field has a direction associated with it, such
that the field exits the North end of a magnet, flows through the air or other materials
nearby, and re-enters the South end of the magnet. Inside the magnet, the field flows
from the South back to the North.
In summary, current flows from + to - of a battery, and magnetic fields flow from the
North to the South of a magnet.
Saturation
Is there a limit to magnetic field strength? What happens to the projectile when the coil
is boosted to an extremely strong magnetic field?

Are you familiar with B-H curves? First, a quick refresher about the definitions of B and
H.
The magnetic field, H, tends to magnetize space.
The magnetic flux density, B, is the total magnetic effect that results.
The two concepts are approximately linked with a simple equation. The magnetic flux
density, B, is the total magnetic effect that results in iron: B = u 0 * ur * H, where u0 is the
permeability of free space, and ur is the relative permeability of the material.
The constant ur is essentially unity (= 1.0) for all materials except ferromagnetic
materials, notably iron, for which the value of the relative permeability lies typically in
the range 1000 to 10,000. Thus for iron, a small magnetic cause, H, creates a large
magnetic effect, B.
By the way, an iron/nickel compound typically has twice the constant relative
permeability (ur) as mere iron. The highest ur is "Supermalloy" at about 100,000 (see
table of magnetic materials) So you might think it would perform much better than a
plain iron projectile. Unfortunately, it ain't necessarily so. The total flux in a magnetic
circuit is limited by any air gaps, and we've got a big one: right down the middle of the
coil.
But a more complete picture takes into account saturation. As the applied field increases,
more and more magnetic domains within the ferrous material will align themselves with
the external field. This temporarily makes the core act as a magnet itself, which for us
will improve the forces upon the projectile. At small to medium applied fields, the flux
density B is nearly linear with respect to the applied field H. But at very strong applied
fields, almost all the magnetic domains are aligned, and it gets much harder to coerce the
last few domains into alignment.
The B-H curve here illustrates the effect of magnetic saturation. It shows the effect of
applying an external magnetic field to unmagnetized iron. The magnetism curve starts at
the origin (0,0) and increases linearly as the magnetic field magnetizes the iron. The
slope of the curve in this region is u = u 0 ur. Eventually, however, all the magnetic

domains align with the applied field, and the curve flattens out as the iron becomes
magnetically saturated. As we reverse the applied magnetic field, the iron eventually
becomes magnetized in the reverse direction until it again saturates.
How does this relate to our projectiles? It means measurements of forces taken at small
coil currents are linear with respect to applied current. As you increase current, the
magnetizing force produced by an air coil is always linear, but the response in the
projectile is not linear. At very large fields, the force in the projectile is not as great as
linear equations would lead you to expect. The mechanical force of attraction does
continue to increase, but it takes more and more coil current to achieve the same amount
of increase.
Conclusion
So how do you apply that general concept to a physical system? The magnetic saturation
point puts a limit on the maximum effective strength to drive your coil. There will be no
benefit at all if you exceed this, only wasted heat.
Basically it goes like this. You need to know the exact material and dimensions of your
projectile. From that, manufacturers supply the ferromagnetic permeability (or see the
table of magnetic materials) and the B-H curve. From the saturation point, Bs, you know
the upper limit for your design; it would be wasteful to attempt any higher flux density.
You also need to compute the total reluctance R along the path of your magnetic circuit.
From Bs and R you can compute the maximum flux your coil should generate. But don't
worry: this upper limit is going to be huge! My initial calculation found it will be
millions of ampere-turns, using any air-core coil that we've seen so far!
Hystersis
Ferromagnetic and ferrimagnetic materials have non-linear initial magnetisation curves
(i.e. the dotted lines in figure 7), as the changing magnetisation with applied field is due
to a change in the magnetic domain structure. These materials also show hysteresis and
the magnetisation does not return to zero after the application of a magnetic field. Figure
7 shows a typical hysteresis loop; the two loops represent the same data, however, the

red loop is the polarisation (J = oM = B - oH) and the blue loop the induction, both
plotted against the applied field.
Figure 7: A typical hysteresis loop for a ferro- or ferrimagnetic material.

Illustrated in the first quadrant of the loop is the initial magnetisation curve (dotted line),
which shows the increase in polarisation (and induction) on the application of a field to
an unmagnetised sample. In the first quadrant the polarisation and applied field are both
positive, i.e. they are in the same direction. The polarisation increases initially by the
growth of favourably oriented domains, which will be magnetised in the easy direction
of the crystal. When the polarisation can increase no further by the growth of domains,
the direction of magnetisation of the domains then rotates away from the easy axis to
align with the field. When all of the domains have fully aligned with the applied field
saturation is reached and the polarisation can increase no further.
If the field is removed the polarisation returns along the solid red line to the y-axis (i.e.
H=0), and the domains will return to their easy direction of magnetisation, resulting in a
decrease in polarisation. In figure 7, the line from the saturation point to the y-axis is

horizontal, which is representative of a well aligned material, where the domains are
magnetised in the easy direction of the crystal at the saturation point.
If the direction of applied field is reversed (i.e. to the negative direction) then the
polarisation will follow the red line into the second quadrant. The hysteresis means that
the polarisation lags behind the applied field and will not immediately switch direction
into the third quadrant (i.e. negative polarisation). The polarisation will only decrease
after a sufficiently high field is applied to:
1. nucleate and grow domains favourably oriented with respect to the applied field or
2. rotate the direction of magnetisation of the domains towards the applied field
After applying a high enough field saturation polarisation will be achieved in the
negative direction. If the applied field is then decreased and again applied in the positive
direction then the full hysteresis loop is plotted.
If the field is repeatedly switched from positive to negative directions and is of sufficient
magnitude then the polarisation and induction will cycle around the hysteresis loop in an
anti-clockwise direction. The area contained within the loop indicates the amount of
energy absorbed by the material during each cycle of the hysteresis loop.
When a ferromagnetic material is magnetized in one direction, it will not relax back to
zero magnetization when the imposed magnetizing field is removed. It must be driven
back to zero by a field in the opposite direction. If an alternating magnetic field is
applied to the material, its magnetization will trace out a loop called a hysteresis loop.
The lack of retraceability of the magnetization curve is the property called hysteresis and
it is related to the existence of magnetic domains in the material. Once the magnetic
domains are reoriented, it takes some energy to turn them back again. This property of
ferrromagnetic materials is useful as a magnetic "memory". Some compositions of
ferromagnetic materials will retain an imposed magnetization indefinitely and are useful
as "permanent magnets". The magnetic memory aspects of iron and chromium oxides

make them useful in audio tape recording and for the magnetic storage of data on
computer disks.
Remanent
So, as we've seen, if we have a magnetic material and place it in an external magnetic
field (one that we've called the inducing field), we can make the magnetic material
produce its own magnetic field. If we were to measure the total magnetic field near the
material, that field would be the sum of the external, or inducing field, and the induced
field produced in the material. By measuring spatial variations in the total magnetic field
and by knowing what the inducing field looks like, we can, in principle, map spatial
variations in the induced field and from this determine spatial variations in the magnetic
susceptibility of the subsurface.
Although this situation is a bit more complex than the gravitational situation, it's still
manageable. There is, however, one more complication in nature concerning material
magnetism that we need to consider. In the scenario we've been discussing, the induced
magnetic field is a direct consequence of a magnetic material being surrounded by an
inducing magnetic field. If you turn off the inducing magnetic field, the induced
magnetization disappears. Or does it?
If the magnetic material has relatively large susceptibilities, or if the inducing field is
strong, the magnetic material will retain a portion of its induced magnetization even after
the induced field disappears. This remaining magnetization is called Remanent
Magnetization.
Remanent Magnetization is the component of the material's magnetization that solid-

earth geophysicists use to map the motion of continents and ocean basins resulting from
plate tectonics. Rocks can acquire a remanent magnetization through a variety of
processes that we don't need to discuss in detail. A simple example, however, will
illustrate the concept. As a volcanic rock cools, its temperature decreases past the Curie

Temperature. At the Curie Temperature, the rock, being magnetic, begins to produce an
induced magnetic field. In this case, the inducing field is the Earth's magnetic field. As
the Earth's magnetic field changes with time, a portion of the induced field in the rock
does not change but remains fixed in a direction and strength reflective of the Earth's
magnetic field at the time the rock cooled through its Curie Temperature. This is the
remanent magnetization of the rock - the recorded magnetic field of the Earth at the time
the rock cooled past its Curie Temperature.
The only way you can measure the remanent magnetic component of a rock is to take a
sample of the rock back to the laboratory for analysis. This is time consuming and
expensive. As a result, in exploration geophysics, we typically assume there is no
remanent magnetic component in the observed magnetic field. Clearly, however, this
assumption is wrong and could possibly bias our interpretations
The magnetization retained by rocks from previous magnetic fields,abbreviated NRM.

NRM is a record of the Earth's magnetic field as it existed at the time that the rock
formed, such as when magnetic crystals in igneous rocks solidified (also known as
chemical remanent magnetism, CRM) or at the time of deposition of sedimentary rocks
(known as depositional remanent magnetism, DRM). During deposition of sediments
that become sedimentary rock, magnetized particles can settle with their magnetic pole
aligned with that of the Earth at that time.
Magnetic Flux Density

Magnetic-flux-density sensors provide an electrical output proportional to the flux
density. The most popular device for measuring the flux density is the Hall-effect sensor,
in which the excitation current is held constant and the semiconductor crystal placed in
the field generates voltage along the axis perpendicular to the direction of the field and
current flow. This voltage is proportional to the flux density. The devices are designed as
small-size probes, which can contain in their tip one, two, or three crystals for measuring
the field in one, two, or three mutually orthogonal transverse directions. An inductor
probe is made of an air-core inductor. The voltage induced in a coil placed in an

alternating magnetic field is proportional to the measured flux density. Steady-state

magnetic fields can be measured by spinning the coil in the field at a constant speed. The
ac voltage developed in the coil due to crossing the lines of the field is a function of the
flux density. A nuclear magnetic resonance magnetometer is based on the effect of the
sensitivity of the nuclear resonance frequency of certain substances placed in the
magnetic field to the strength of the field. The transduction element consists of a small
container with a substance in it such as deuterium or lithium. A coil, excited by a radio-
frequency current, is wound around the container. The resonant frequency is identified
by a sharp increase in the power consumed by the substance because of the characteristic
energy absorption at resonance. Counting the frequency allows measurement of the flux
density. A flux-gate magnetometer, used for measuring extremely small flux densities, is
composed of several saturable reactors that are excited by an ac voltage. This voltage
keeps the induction in the cores of the reactors close to saturation. When the outside
steady field acts against the core, the component of a second-harmonic current is
produced in the reactor's circuit. It is noteworthy that this component does not manifest
itself at all in the absence of the outside field. With the increase of the field intensity, the
second-harmonic voltage can be detected, giving a measure of the flux density Three
mutually orthogonal reactors allow measurement of the flux density along the three
special axes.
Coercive Power
Coercion can take many forms. I may prevent you from doing something you wish to do,
by withholding some resources or by physically constraining you. For example, the
modern state imprisons those who do not act in accordance with its legal mandates. In
other cases, I may push you into a behavior in which you would otherwise not engage.
For example, parents may use a variety of strategies for getting a resistant child to go to
school, including physically taking the child to the school building. As another example,
the majority of nations of the world joined in a boycott of Iraqi oil, in the hope of forcing
the Iraqi government to honor the peace agreement that ended the 1991 Gulf War.

While not all of these forms are typically categorized as violent, coercion is usually
associated with physical violence. As C. Wright Mills says, "All politics is a struggle for
power; the ultimate kind of power is violence."[7] Violence can produce changes in the
target. The slave who is whipped may return to work, at least make the attempt to show
compliance while the overseer is watching, and try to avoid additional lashes. A prisoner
who is tortured may divulge sought-after information in order to end the torture. The
warring enemy may sign a truce, because it no longer has resources to continue the
fighting.
Coercive power is most effective, however, when the threat of violence or other
punishment is sufficient in itself to get the target to accede to the demand
Coercivity and Remanence in Permanent Magnets

A good permanent magnet should produce a high magnetic field with a low mass, and
should be stable against the influences which would demagnetize it. The desirable
properties of such magnets are typically stated in terms of the remanence and coercivity
of the magnet materials.
When a ferromagnetic material is magnetized in one direction, it will not relax back to
zero magnetization when the imposed magnetizing field is removed. The amount of
magnetization it retains at zero driving field is called its remanence. It must be driven
back to zero by a field in the opposite direction; the amount of reverse driving field
required to demagnetize it is called its coercivity. If an alternating magnetic field is

applied to the material, its magnetization will trace out a loop called a hysteresis loop.
The lack of retraceability of the magnetization curve is the property called hysteresis and
it is related to the existence of magnetic domains in the material. Once the magnetic
domains are reoriented, it takes some energy to turn them back again. This property of
ferrromagnetic materials is useful as a magnetic "memory". Some compositions of
ferromagnetic materials will retain an imposed magnetization indefinitely and are useful
as "permanent magnets".
The table below contains some data about materials used as permanent magnets. Both
the coercivity and remanence are quoted in Tesla, the basic unit for magnetic field B.
The hysteresis loop above is plotted in the form of magnetization M as a function of
driving magnetic field strength H. This practice is commonly followed because it shows
the external driving influence (H) on the horizontal axis and the response of the material
(M) on the vertical axis. Besides coercivity and remanence, a quality factor for
permanent magnets is the quantity (BB0/μ0)max. A high value for this quantity implies that
the required magnetic flux can be obtained with a smaller volume of the material,
making the device lighter and more compact.
Coercivity Remanence (BB0/μ0)max

Material
(T) (T) (kJ/m3)
BaFe12O19 0.36 0.36 25
Alnico IV 0.07 0.6 10.3
Alnico V 0.07 1.35 55
Alcomax I 0.05 1.2 27.8
MnBi 0.37 0.48 44
Ce(CuCo)5 0.45 0.7 92
SmCo5 1.0 0.83 160
Sm2Co17 0.6 1.15 215
Nd2Fe14B 1.2 1.2 260

Data from Myers

The alloys from which permanent magnets are made are often very difficult to handle
metallurgically. They are mechanically hard and brittle. They may be cast and then
ground into shape, or even ground to a powder and formed. From powders, they may be
mixed with resin binders and then compressed and heat treated. Maximum anisotropy of
the material is desirable, so to that end the materials are often heat treated in the presence
of a strong magnetic field.
The materials with high remanence and high coercivity from which permanent magnets
are made are sometimes said to be "magnetically hard" to contrast them with the
"magnetically soft" materials from which transformer cores and coils for electronics are
made.

Applied - Physics - Notes 1

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The Great Eastern Shipping Co. Ltd.

Forces; Newton's second law

where is the acceleration. It is not always the case that m is independent of t.

where Ek is called the kinetic energy. For a point particle, it is defined as

Developed by the ISF Maritime Services Pvt. Ltd -1-

A particular class of forces, known as conservative forces, can be expressed as the

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joule J = N·m m2·kg·s−2 Energy

Orbital Energies, Kepler's Laws and Other Relationships

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Movement of a planet or satellite in an

The semi-major axis is directly related to the

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Keplers 2nd Law

Kepler's 3rd Law - Relationship between Period and Semi-major Axis

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The idea of such a "space tether"?

--The calculation is very similar to the preceding one.

We have 43200 sec = 5063.5 R*SQRT(R) = 5063.5 R3/2

R3/2 = 8.53164 R = (8.53164)2/3 = 4.175 RE

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The constant of proportionality G is known as the universal gravitational constant. It is

Newton's Modification of Kepler's Third Law

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Weight and the Gravitational Force

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perpendicular distance) and is the Kronecker delta, or

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(5)which can be written as

(7)where is the angular velocity vector.

(9)which is an eigenvalue problem.

cylinder about central diameter

ellipsoid about principal axis

elliptical slab about major axis

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elliptical slab about vertical

rectangular parallelepiped about major axis

ring about perpendicular axis

ring about diameter

rod about end

rod about center

sphere about diameter

torus about diameter

torus about symmetry axis

rectangle about minor axis

Developed by the ISF Maritime Services Pvt. Ltd - 12 -

Moment of Inertia and Angular Momentum

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axis on the far end of a thin rod:

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solid disks and cylinders

Both objects have the same mass and equal diameters.

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For a point mass, its angular momentum can be expressed as

L =  where = mr2 and  = v/r

Law of Conservation of Angular Momentum

Angular momentum is conserved whenever there is no external force exerting a torque

Conservation of angular momentum justifies the relationship shown in Kepler's Second

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The radius of gyration (r)

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Regardless of the magnitude of the force and displacement, Fdcosine 90 degrees is 0