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You are on page 1of 1290

1)

– volume 2

Chapter 242

**16-00 – General reference works
**

(handbooks, dictionaries,

bibliographies, etc.)

242.1 direct product of modules

**Let {Xi : i ∈ Q I} be a collection of modules in some category of modules. Then the
**

direct product i∈I Xi of that collection is the module whose underlying set is the cartesian product

of the Xi with componentwise addition and scalar multiplication. For example, in a category

of left modules:

(xi ) + (yi ) = (xi + yi),

r(xi ) = (rxi ).

Q

For each j ∈QI we have a projection pj : i∈I Xi → Xj defined by (xi ) 7→ xQ

j , and an injection

λj : Xj → i∈I Xi where an element xj of Xj maps to the element of i∈I Xi whose jth

term is xj and every other term is zero.

Q

The direct product i∈I Xi satisfies a certain universal property. Namely, if Y is a module

and there exist homomorphisms

Q fi : Xi → Y for all i ∈ I, then there exists a unique

homomorphism φ : Y → i∈I Xi satisfying φλi = fi for all i ∈ I.

fi

Xi Y

λi φ

Q

i∈I Xi

**The direct product is often referred to as the complete direct sum, or the strong direct sum,
**

or simply the product.

1088

Compare this to the direct sum of modules.

Version: 3 Owner: antizeus Author(s): antizeus

242.2 direct sum

**Let {X` i : i ∈ I} be a collection of modules in some category of modules. Then the direct
**

sum i∈I Xi of that collection is the submodule of the direct product of the Xi consisting

of all elements (xi ) such that all but a finite number of the xi are zero.

`

For each j ∈`I we have a projection pj : i∈I Xi → Xj defined by (xi ) 7→ x` j , and an injection

λj : Xj → i∈I Xi where an element xj of Xj maps to the element of i∈I Xi whose jth

term is xj and every other term is zero.

`

The direct sum i∈I Xi satisfies a certain universal property. Namely, if Y is a module

and there exist homomorphisms

` fi : Y → Xi for all i ∈ I, then there exists a unique

homomorphism φ : i∈I Xi → Y satisfying pi φ = fi for all i ∈ I.

fi

Xi Y

pi φ

`

i∈I Xi

The direct sum is often referred to as the weak direct sum or simply the sum.

Compare this to the direct product of modules.

Version: 3 Owner: antizeus Author(s): antizeus

242.3 exact sequence

**If we have two homomorphisms f : A → B and g : B → C in some category of modules,
**

then we say that f and g are exact at B if the image of f is equal to the kernel of g.

A sequence of homomorphisms

fn+1 fn

· · · → An+1 −→ An −→ An−1 → · · ·

is said to be exact if each pair of adjacent homomorphisms (fn+1 , fn ) is exact – in other

words if imfn+1 = kerfn for all n.

Compare this to the notion of a chain complex.

Version: 2 Owner: antizeus Author(s): antizeus

1089

242.4 quotient ring

Definition.

Let R be a ring and let I be a two-sided ideal of R. To define the quotient ring R/I, let us first

define an equivalence relation in R. We say that the elements a, b ∈ R are equivalent, written

as a ∼ b, if and only if a − b ∈ I. If a is an element of R, we denote the corresponding

equivalence class by [a]. Thus [a] = [b] if and only if a − b ∈ I. The quotient ring of R

modulo I is the set R/I = {[a] | a ∈ R}, with a ring structure defined as follows. If [a], [b]

are equivalence classes in R/I, then

• [a] + [b] := [a + b],

• [a] · [b] := [a · b].

**Here a and b are some elements in R that represent [a] and [b]. By construction, every
**

element in R/I has such a representative in R. Moreover, since I is closed under addition

and multiplication, one can verify that the ring structure in R/I is well defined.

properties.

1. If R is commutative, then R/I is commutative.

Examples.

1. For any ring R, we have that R/R = {0} and R\{0} = R.

**2. Let R = Z, and let I be the set of even numbers. Then R/I contains only two classes;
**

one for even numbers, and one for odd numbers.

Version: 3 Owner: matte Author(s): matte, djao

1090

Chapter 243

16D10 – General module theory

243.1 annihilator

Let R be a ring.

Suppose that M is a left R-module.

**If X is a subset of M, then we define the left annihilator of X in R:
**

l.ann(X) = {r ∈ R | rx = 0 for all x ∈ X}.

**If Z is a subset of R, then we define the right annihilator of Z in M:
**

r.annM (Z) = {m ∈ M | zm = 0 for all z ∈ Z}.

Suppose that N is a right R-module.

**If Y is a subset of N, then we define the right annihilator of Y in R:
**

r.ann(Y ) = {r ∈ R | yr = 0 for all y ∈ Y }.

**If Z is a subset of R, then we define the left annihilator of Z in N:
**

l.annN (Z) = {n ∈ N | nz = 0 for all z ∈ Z}.

Version: 3 Owner: antizeus Author(s): antizeus

243.2 annihilator is an ideal

The right annihilator of a right R-module MR in R is an ideal.

1091

B y the distributive law for modules, it is easy to see that r. ann(MR ) is closed under

addition and right multiplication. Now take x ∈ r. ann(MR ) and r ∈ R.

**Take any m ∈ MR . Then mr ∈ MR , but then (mr)x = 0 since x ∈ r. ann(MR ). So m(rx) = 0
**

and rx ∈ r. ann(MR ).

An equivalent result holds for left annihilators.

Version: 2 Owner: saforres Author(s): saforres

243.3 artinian

A module M is artinian if it satisfies the following equivalent conditions:

**• the descending chain condition holds for submodules of M;
**

• every nonempty family of submodules of M has a minimal element.

**A ring R is left artinian if it is artinian as a left module over itself (i.e. if R R is an artinian
**

module), and right artinian if it is artinian as a right module over itself (i.e. if RR is an

artinian module), and simply artinian if both conditions hold.

Version: 3 Owner: antizeus Author(s): antizeus

243.4 composition series

**Let R be a ring and let M be a (right or left) R-module. A series of submodules
**

M = M0 ⊃ M1 ⊃ M2 ⊃ · · · ⊃ Mn = 0

in which each quotient Mi /Mi+1 is simple is called a composition series for M.

**A module need not have a composition series. For example, the ring of integers, Z, cond-
**

sidered as a module over itself, does not have a composition series.

**A necessary and sufficient condition for a module to have a composition series is that it is
**

both noetherian and artinian.

If a module does have a composition series, then all composition series are the same length.

This length (the number n above) is called the composition length of the module.

If R is a semisimple Artinian ring, then RR and R R always have composition series.

Version: 1 Owner: mclase Author(s): mclase

1092

243.5 conjugate module

**If M is a right module over a ring R, and α is an endomorphism of R, we define the
**

conjugate module M α to be the right R-module whose underlying set is {mα | m ∈ M},

with abelian group structure identical to that of M (i.e. (m − n)α = mα − nα ), and scalar

multiplication given by mα · r = (m · α(r))α for all m in M and r in R.

**In other words, if φ : R → EndZ (M) is the ring homomorphism that describes the right module action
**

of R upon M, then φα describes the right module action of R upon M α .

If N is a left R-module, we define α N similarly, with r · α n = α (α(r) · n).

Version: 4 Owner: antizeus Author(s): antizeus

243.6 modular law

**Let R M be a left R-module with submodules A, B, C, and suppose C ⊆ B. Then
**

\ \

C + (B A) = B (C + A)

Version: 1 Owner: saforres Author(s): saforres

243.7 module

Let R be a ring, and let M be an abelian group.

**We say that M is a left R-module if there exists a ring homomorphism φ : R → EndZ (M)
**

from R to the ring of abelian group endomorphisms on M (in which multiplication of endo-

morphisms is composition, using left function notation). We typically denote this function

using a multiplication notation:

[φ(r)](m) = r · m = rm

This ring homomorphism defines what is called a left module action of R upon M.

**If R is a unital ring (i.e. a ring with identity), then we typically demand that the ring
**

homomorphism map the unit 1 ∈ R to the identity endomorphism on M, so that 1 · m = m

for all m ∈ M. In this case we may say that the module is unital.

**Typically the abelian group structure on M is expressed in additive terms, i.e. with operator
**

+, identity element 0M (or just 0), and inverses written in the form −m for m ∈ M.

1093

Right module actions are defined similarly, only with the elements of R being written on

the right sides of elements of M. In this case we either need to use an anti-homomorphism

R → EndZ (M), or switch to right notation for writing functions.

Version: 7 Owner: antizeus Author(s): antizeus

**243.8 proof of modular law
**

T T

T A) ⊆ B (C + A):

First we show C + (B T

Note that C ⊆ B, B AT⊆ B, and therefore C + (BT A) ⊆ B.

Further, C ⊆ C + A, B A ⊆ C + A, thus C + (B A) ⊆ C + A.

T T

Next we show

T B (C + A) ⊆ C + (B A):

Let b ∈ B (C + A). Then b = c + a for some c ∈ C and a ∈ A. Hence a = b − c, and so

a ∈ B since bT∈ B and c ∈ C ⊆ B. T

Hence a ∈ B A, so b = c + a ∈ C + (B A).

Version: 5 Owner: saforres Author(s): saforres

243.9 zero module

Let R be a ring.

**The abelian group which contains only an identity element (zero) gains a trivial R-module
**

structure, which we call the zero module.

**Every R-module M has an zero element and thus a submodule consisting of that element.
**

This is called the zero submodule of M.

Version: 2 Owner: antizeus Author(s): antizeus

1094

Chapter 244

16D20 – Bimodules

244.1 bimodule

**Suppose that R and S are rings. An (R, S)-bimodule is an abelian group M which has a left
**

R-module action as well as a right S-module action, which satisfy the relation r(ms) = (rm)s

for every choice of elements r of R, s of S, and m of M.

**A (R, S)-sub-bi-module of M is a subgroup which is also a left R-submodule and a right
**

S-submodule.

Version: 3 Owner: antizeus Author(s): antizeus

1095

Chapter 245

16D25 – Ideals

245.1 associated prime

**Let R be a ring, and let M be an R-module. A prime ideal P of R is an annihilator prime
**

for M if P is equal to the annihilator of some nonzero submodule X of M.

**Note that if this is the case, then the module annA (P ) contains X, has P as its annihilator,
**

and is a faithful (R/P )-module.

**If, in addition, P is equal to the annihilator of a submodule of M that is a fully faithful
**

(R/P )-module, then we call P an associated prime of M.

Version: 2 Owner: antizeus Author(s): antizeus

245.2 nilpotent ideal

**A left (right) ideal I of a ring R is a nilpotent ideal if I n = 0 for some positive integer n.
**

Here I n denotes a product of ideals – I · I · · · I.

Version: 2 Owner: antizeus Author(s): antizeus

245.3 primitive ideal

**Let R be a ring, and let I be an ideal of R. We say that I is a left (right) primitive ideal if
**

there exists a simple left (right) R-module X such that I is the annihilator of X in R.

We say that R is a left (right) primitive ring if the zero ideal is a left (right) primitive ideal

1096

of R.

Note that I is a left (right) primitive ideal if and only if R/I is a left (right) primitive ring.

Version: 2 Owner: antizeus Author(s): antizeus

245.4 product of ideals

**Let R be a ring, and let A and B be left (right) ideals of R. Then the product of the
**

ideals A and B, which we denote AB, is the left (right) ideal generated by the products

{ab | a ∈ A, b ∈ B}.

Version: 2 Owner: antizeus Author(s): antizeus

245.5 proper ideal

**Suppose R is a ring and I is an ideal of R. We say that I is a proper ideal if I is not equal
**

to R.

Version: 2 Owner: antizeus Author(s): antizeus

245.6 semiprime ideal

**Let R be a ring. An ideal I of R is a semiprime ideal if it satisfies the following equivalent
**

conditions:

(a) I can be expressed as an intersection of prime ideals of R;

(b) if x ∈ R, and xRx ⊂ I, then x ∈ I;

(c) if J is a two-sided ideal of R and J 2 ⊂ I, then J ⊂ I as well;

(d) if J is a left ideal of R and J 2 ⊂ I, then J ⊂ I as well;

(e) if J is a right ideal of R and J 2 ⊂ I, then J ⊂ I as well.

Here J 2 is the product of ideals J · J.

**The ring R itself satisfies all of these conditions (including being expressed as an intersection
**

of an empty family of prime ideals) and is thus semiprime.

A ring R is said to be a semiprime ring if its zero ideal is a semiprime ideal.

1097

Note that an ideal I of R is semiprime if and only if the quotient ring R/I is a semiprime

ring.

Version: 7 Owner: antizeus Author(s): antizeus

245.7 zero ideal

In any ring, the set consisting only of the zero element (i.e. the additive identity) is an ideal

of the left, right, and two-sided varieties. It is the smallest ideal in any ring.

Version: 2 Owner: antizeus Author(s): antizeus

1098

Chapter 246

**16D40 – Free, projective, and flat
**

modules and ideals

246.1 finitely generated projective module

**Let R be a unital ring. A finitely generated projective right R-module is of the form eRn ,
**

n ∈ N, where e is an idempotent in EndR (Rn ).

**Let A be a unital C ∗ -algebra and p be a projection in EndA (An ), n ∈ N. Then, E = pAn is
**

a finitely generated projective right A-module. Further, E is a pre-Hilbert A-module with

(A-valued) inner product

n

X

hu, vi = u∗i vi , u, v ∈ E.

i=1

Version: 3 Owner: mhale Author(s): mhale

246.2 flat module

**A right module M over a ring R is flat if the tensor product functor M ⊗R (−) is an
**

exact functor.

**Similarly, a left module N over R is flat if the tensor product functor (−) ⊗R N is an exact
**

functor.

Version: 2 Owner: antizeus Author(s): antizeus

1099

246.3 free module

**Let R be a commutative ring. A free module over R is a direct sum of copies of R. In
**

particular, as every abelian group is a Z-module, a free abelian group is a direct sum of

copies of Z. This is equivalent to saying that the module has a free basis, i.e. a set of

elements with the property that every element of the module can be uniquely expressed as

an linear combination over R of elements of the free basis. In the case that a free module

over R is a sum of finitely many copies of R, then the number of copies is called the rank of

the free module.

**An alternative definition of a free module is via its universal property: Given a set X, the
**

free R-module F (X) on the set X is equipped with a function i : X → F (X) satisfying the

property that for any other R-module A and any function f : X → A, there exists a unique

R-module map h : F (X) → A such that (h ◦ i) = f .

Version: 4 Owner: mathcam Author(s): mathcam, antizeus

246.4 free module

Let R be a ring. A free module over R is a direct sum of copies of R.

**Similarly, as an abelian group is simply a module over Z, a free abelian group is a direct sum
**

of copies of Z.

**This is equivalent to saying that the module has a free basis, i.e. a set of elements with the
**

property that every element of the module can be uniquely expressed as an linear combination

over R of elements of the free basis.

Version: 1 Owner: antizeus Author(s): antizeus

246.5 projective cover

**Let X and P be modules. We say that P is a projective cover of X if P is a projective module
**

and there exists an epimorphism p : P → X such that ker p is a superfluous submodule of P .

**Equivalently, P is an projective cover of X if P is projective, and there is an epimorphism
**

p : P → X, and if g : P 0 → X is an epimorphism from a projective module P 0 to X, then

1100

there exists an epimorphism h : P 0 → P such that ph = g.

P0

h g

P p X 0

0

Version: 2 Owner: antizeus Author(s): antizeus

246.6 projective module

A module P is projective if it satisfies the following equivalent conditions:

(a) Every short exact sequence of the form 0 → A → B → P → 0 is split;

(b) The functor Hom(P, −) is exact;

**(c) If f : X → Y is an epimorphism and there exists a homomorphism g : P → Y , then
**

there exists a homomorphism h : P → X such that f h = g.

P

h g

X f

Y 0

(d) The module P is a direct summand of a free module.

Version: 3 Owner: antizeus Author(s): antizeus

1101

Chapter 247

**16D50 – Injective modules,
**

self-injective rings

247.1 injective hull

**Let X and Q be modules. We say that Q is an injective hull or injective envelope of X if Q
**

is both an injective module and an essential extension of X.

**Equivalently, Q is an injective hull of X if Q is injective, and X is a submodule of Q, and
**

if g : X → Q0 is a monomorphism from X to an injective module Q0 , then there exists a

monomorphism h : Q → Q0 such that h(x) = g(x) for all x ∈ X.

0

i

0 X Q

g

h

0

Q

Version: 2 Owner: antizeus Author(s): antizeus

247.2 injective module

A module Q is injective if it satisfies the following equivalent conditions:

(a) Every short exact sequence of the form 0 → Q → B → C → 0 is split;

(b) The functor Hom(−, Q) is exact;

1102

(c) If f : X → Y is a monomorphism and there exists a homomorphism g : X → Q, then

there exists a homomorphism h : Y → Q such that hf = g.

f

0 X Y

g

h

Q

Version: 3 Owner: antizeus Author(s): antizeus

1103

Chapter 248

**16D60 – Simple and semisimple
**

modules, primitive rings and ideals

248.1 central simple algebra

**Let K be a field. A central simple algebra A (over K) is an algebra A over K, which is
**

finite dimensional as a vector space over K, such that

• A has an identity element, as a ring

**• A is central: the center of A equals K (for all z ∈ A, we have z · a = a · z for all a ∈ A
**

if and only if z ∈ K)

• A is simple: for any two sided ideal I of A, either I = {0} or I = A

**By a theorem of Brauer, for every central simple algebra A over K, there exists a unique (up
**

to isomorphism) division ring D containing K and a unique natural number n such that A

is isomorphic to the ring of n × n matrices with coefficients in D.

Version: 2 Owner: djao Author(s): djao

248.2 completely reducible

**A module M is called completely reducible (or semisimple) if it is a direct sum of irreducible
**

(or simple) modules.

Version: 1 Owner: bwebste Author(s): bwebste

1104

248.3 simple ring

**A nonzero ring R is said to be a simple ring if it has no (two-sided) ideal other then the
**

zero ideal and R itself.

This is equivalent to saying that the zero ideal is a maximal ideal.

If R is a commutative ring with unit, then this is equivalent to being a field.

Version: 4 Owner: antizeus Author(s): antizeus

1105

Chapter 249

**16D80 – Other classes of modules and
**

ideals

249.1 essential submodule

**Let X be a submodule of a module Y . We say that X is an essential submodule of Y , and
**

that

T Y is an essential extension of X, if whenever A is a nonzero submodule of Y , then

A X is also nonzero.

**A monomorphism f : X → Y is an essential monomorphism if the image imf is an essential
**

submodule of Y .

Version: 2 Owner: antizeus Author(s): antizeus

249.2 faithful module

**Let R be a ring, and let M be an R-module. We say that M is a faithful R-module if its
**

annihilator annR (M) is the zero ideal.

We say that M is a fully faithful R-module if every nonzero R-submodule of M is faithful.

Version: 3 Owner: antizeus Author(s): antizeus

1106

249.3 minimal prime ideal

**A prime ideal P of a ring R is called a minimal prime ideal if it does not properly contain
**

any other prime ideal of R.

If R is a prime ring, then the zero ideal is a prime ideal, and is thus the unique minimal

prime ideal of R.

Version: 2 Owner: antizeus Author(s): antizeus

249.4 module of finite rank

**Let M be a module, and let E(M) be the injective hull of M. Then we say that M has finite
**

rank if E(M) is a finite direct sum of indecomposible submodules.

**This turns out to be equivalent to the property that M has no infinite direct sums of nonzero
**

submodules.

Version: 3 Owner: antizeus Author(s): antizeus

249.5 simple module

**Let R be a ring, and let M be an R-module. We say that M is a simple or irreducible module
**

if it contains no submodules other than itself and the zero module.

Version: 2 Owner: antizeus Author(s): antizeus

249.6 superfluous submodule

**Let X be a submodule of a module Y . We say that X is a superfluous submodule of Y if
**

whenever A is a submodule of Y such that A + X = Y , then A = Y .

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249.7 uniform module

**A module M is said to be uniform if any two nonzero submodules of M must have a nonzero
**

intersection. This is equivalent to saying that any nonzero submodule is an essential submodule.

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1108

Chapter 250

**16E05 – Syzygies, resolutions,
**

complexes

250.1 n-chain

**An n-chain on a topological space X is a finite formal sum of n-simplices on X. The group
**

of such chains is denoted Cn (X). For a CW-complex Y, Cn (Y ) = Hn (Y n , Y n−1 ), where Hn

denotes the nth homology group.

**The boundary of an n-chain is the (n − 1)-chain given by the formal sum of the boundaries
**

of its constitutent simplices. An n-chain is closed if its boundary is 0 and exact if it is the

boundary of some (n + 1)-chain.

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250.2 chain complex

**A sequence of modules and homomorphisms
**

dn+1 n d

· · · → An+1 −→ An −→ An−1 → · · ·

**is said to be a chain complex or complex if each pair of adjacent homomorphisms (dn+1 , dn )
**

satisfies the relation dn dn+1 = 0. This is equivalent to saying that im dn+1 ⊂ ker dn . We

often denote such a complex as (A, d) or simply A.

Compare this to the notion of an exact sequence.

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250.3 flat resolution

Let M be a module. A flat resolution of M is an exact sequence of the form

· · · → Fn → Fn−1 → · · · → F1 → F0 → M → 0

where each Fn is a flat module.

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250.4 free resolution

Let M be a module. A free resolution of M is an exact sequence of the form

· · · → Fn → Fn−1 → · · · → F1 → F0 → M → 0

where each Fn is a free module.

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250.5 injective resolution

Let M be a module. An injective resolution of M is an exact sequence of the form

0 → M → Q0 → Q1 → · · · → Qn−1 → Qn → · · ·

where each Qn is an injective module.

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250.6 projective resolution

Let M be a module. A projective resolution of M is an exact sequence of the form

· · · → Pn → Pn−1 → · · · → P1 → P0 → M → 0

where each Pn is a projective module.

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250.7 short exact sequence

A short exact sequence is an exact sequence of the form

0 → A → B → C → 0.

**Note that in this case, the homomorphism A → B must be a monomorphism, and the
**

homomorphism B → C must be an epimorphism.

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**250.8 split short exact sequence
**

f g

In an abelian category, a short exact sequence 0 → A → B → C → 0 is split if it satisfies

the following equivalent conditions:

(a) there exists a homomorphism h : C → B such that gh = 1C ;

(b) there exists a homomorphism j : B → A such that jf = 1A ;

(c) B is isomorphic to the direct sum A ⊕ C.

In this case, we say that h and j are backmaps or splitting backmaps.

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250.9 von Neumann regular

**An element a of a ring R is said to be von Neumann regular if there exists b ∈ R such
**

that aba = a.

**A ring R is said to be a von Neumann regular ring (or simply a regular ring, if the
**

meaning is clear from context) if every element of R is von Neumann regular.

**Note that regular ring in the sense of von Neumann should not be confused with regular
**

ring in the sense of commutative algebra.

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1111

Chapter 251

16K20 – Finite-dimensional

251.1 quaternion algebra

**A quaternion algebra over a field K is a central simple algebra over K which is four dimen-
**

sional as a vector space over K.

Examples:

**• For any field K, the ring M2×2 (K) of 2 × 2 matrices with entries in K is a quaternion
**

algebra over K. If K is algebraically closed, then all quaternion algebras over K are

isomorphic to M2×2 (K).

**• For K = R, the well known algebra H of Hamiltonian quaternions is a quaternion
**

algebra over R. The two algebras H and M2×2 (R) are the only quaternion algebras

over R, up to isomorphism.

**• When K is a number field, there are infinitely many non–isomorphic quaternion alge-
**

bras over K. In fact, there is one such quaternion algebra for every even sized finite

collection of finite primes or real primes of K. The proof of this deep fact leads to

many of the major results of class field theory.

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Chapter 252

16K50 – Brauer groups

252.1 Brauer group

**Let K be a field. The Brauer group Br(K) of K is the set of all equivalence classes of
**

central simple algebras over K, where two central simple algebras A and B are equivalent

if there exists a division ring D over K and natural numbers n, m such that A (resp. B) is

isomorphic to the ring of n × n (resp. m × m) matrices with coefficients in D.

**The group operation in Br(K) is given by tensor product: for any two central simple al-
**

gebras A, B over K, their product in Br(K) is the central simple algebra A ⊗K B. The

identity element in Br(K) is the class of K itself, and the inverse of a central simple algebra

A is the opposite algebra Aopp defined by reversing the order of the multiplication operation

of A.

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Chapter 253

16K99 – Miscellaneous

253.1 division ring

A division ring is a ring D with identity such that

• 1 6= 0

• For all nonzero a ∈ D, there exists b ∈ D with a · b = b · a = 1

A field is equivalent to a commutative division ring.

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1114

Chapter 254

**16N20 – Jacobson radical,
**

quasimultiplication

254.1 Jacobson radical

**The Jacobson radical J(R) of a ring R is the intersection of the annihilators of irreducible
**

left R-modules.

The following are alternate characterizations of the Jacobson radical J(R):

1. The intersection of all left primitive ideals.

2. The intersection of all maximal left ideals.

**3. The set of all t ∈ R such that for all r ∈ R, 1 − rt is left invertible (i.e. there exists u
**

such that u(1 − rt) = 1).

4. The largest ideal I such that for all v ∈ I, 1 − v is a unit in R.

5. (1) - (3) with “left” replaced by “right” and rt replaced by tr.

Note that if R is commutative and finitely generated, then

J(R) = {x ∈ R | xn = 0for some n ∈ N} = Nil(R)

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254.2 a ring modulo its Jacobson radical is semiprim-

itive

Let R be a ring. Then J(R/J(R)) = (0).

**L et [u] ∈ J(R/J(R)). Then by one of the alternate characterizations of the Jacobson radical,
**

1 − [r][u] is left invertible for all r ∈ R, so there exists v ∈ R such that [v](1 − [r][u]) = 1.

Then v(1 − ru) = 1 − a for some a ∈ J(R).

**So wv(1−ru) = 1 since w(1−a) = 1 for some w ∈ R. Since this holds for all r ∈ R, u ∈ J(R),
**

then [u] = 0.

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254.3 examples of semiprimitive rings

Examples of semiprimitive rings:

The integers Z:

**S ince Z is commutative, any left ideal is two-sided. So the maximal left ideals
**

T of Z are

the maximal ideals of Z, which are the ideals pZ for p prime. Note that pZ qZ = (0) if

gcd(p, q) > 1. T

Hence J(Z) = p pZ = (0).

A matrix ring Mn (D) over a division ring D:

T he ring Mn (D) is simple, so the only proper ideal is (0). Thus J(Mn (D)) = (0).

A polynomial ring R[x] over a domain R:

**T ake a ∈ J(R[x]) with a 6= 0. Then ax ∈ J(R[x]), since J(R[x]) is an ideal, and deg(ax) >
**

1.

By one of the alternate characterizations of the Jacobson radical, 1 − ax is a unit. But

deg(1 − ax) = max{deg(1), deg(ax)} > 1.

So 1 − ax is not a unit, and by this contradiction we see that J(R[x]) = (0).

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254.4 proof of Characterizations of the Jacobson rad-

ical

**First, note that by definition a left primitive ideal is the annihilator of an irreducible left R-
**

module, so clearly characterization 1) is equivalent to the definition of the Jacobson radical.

**Next, we will prove cyclical containment. Observe that 5) follows after the equivalence of 1)
**

- 4) is established, since 4) is independent of the choice of left or right ideals.

**1) ⊂ 2) We know that every left primitive ideal is the largest ideal contained in a maximal
**

left ideal. So the intersection of all left primitive ideals will be contained in the inter-

section of all maximal left ideals.

T

2) ⊂ 3) Let ST= {M : M a maximal left ideal of R} and take r ∈ R. Let t ∈ M ∈S M. Then

rt ∈ M ∈S M.

Assume 1 − rt is not left invertible; therefore there exists a maximal left ideal M0 of

R such that R(1 − rt) ⊆ M0 .

Note then that 1 − rt ∈ M0 . Also, by definition of t, we have rt ∈ M0 . Therefore

1 ∈ M0 ; this contradiction implies 1 − rt is left invertible.

**3) ⊂ 4) We claim that 3) satisfies the condition of 4).
**

Let K = {t ∈ R : 1 − rt is left invertible for all r ∈ R}.

We shall first show that K is an ideal.

Clearly if t ∈ K, then rt ∈ K. If t1 , t2 ∈ K, then

1 − r(t1 + t2 ) = (1 − rt1 ) − rt2

Now there exists u1 such that u1(1 − rt1 ) = 1, hence

u1 ((1 − rt1 ) − rt2 ) = 1 − u1rt2

Similarly, there exists u2 such that u2 (1 − u1 rt2 ) = 1, therefore

u2 u1 (1 − r(t1 + t2 )) = 1

Hence t1 + t2 ∈ K.

Now if t ∈ K, r ∈ R, to show that tr ∈ K it suffices to show that 1−tr is left invertible.

Suppose u(1 − rt) = 1, hence u − urt = 1, then tur − turtr = tr.

So (1 + tur)(1 − tr) = 1 + tur − tr − turtr = 1.

Therefore K is an ideal.

Now let v ∈ K. Then there exists u such that u(1 − v) = 1, hence 1 − u = −uv ∈ K,

so u = 1 − (1 − u) is left invertible.

So there exists w such that wu = 1, hence wu(1 − v) = w, then 1 − v = w. Thus

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(1 − v)u = 1 and therefore 1 − v is a unit.

Let J be the largest ideal such that, for all v ∈ J, 1 − v is a unit. We claim that

K ⊆ J.

Suppose this were not true; in this case K + J strictly contains J. Consider rx + sy ∈

K + J with x ∈ K, y ∈ J and r, s ∈ R. Now 1 − (rx + sy) = (1 − rx) − sy, and since

rx ∈ K, then 1 − rx = u for some unit u ∈ R.

So 1 − (rx + sy) = u − sy = u(1 − u−1 sy), and clearly u−1sy ∈ J since y ∈ J. Hence

1 − u−1 sy is also a unit, and thus 1 − (rx + sy) is a unit.

Thus 1 − v is a unit for all v ∈ K + J. But this contradicts the assumption that J is

the largest such ideal. So we must have K ⊆ J.

**4) ⊂ 1) We must show that if I is an ideal such that for all u ∈ I, 1 − u is a unit, then
**

I ⊂ ann(R M) for every irreducible left R-module R M.

**Suppose this is not the case, so there exists R M such that I 6⊂ ann(R M). Now we
**

know that ann(R M) is the largest ideal inside some maximal left ideal J of R. Thus

we must also have I 6⊂ J, or else this would contradict the maximality of ann(R M)

inside J.

But since I 6⊂ J, then by maximality I + J = R, hence there exist u ∈ I and v ∈ J

such that u + v = 1. Then v = 1 − u, so v is a unit and J = R. But since J is a proper

left ideal, this is a contradiction.

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254.5 properties of the Jacobson radical

Theorem:

Let R, T be rings and ϕ : R → T be a surjective homomorphism. Then ϕ(J(R)) ⊆ J(T ).

**W e shall use the characterization of the Jacobson radical as the set of all a ∈ R such that
**

for all r ∈ R, 1 − ra is left invertible.

Let a ∈ J(R), t ∈ T . We claim that 1 − tϕ(a) is left invertible:

**Since ϕ is surjective, t = ϕ(r) for some r ∈ R. Since a ∈ J(R), we know 1 − ra is left
**

invertible, so there exists u ∈ R such that u(1 − ra) = 1. Then we have

ϕ(u) (ϕ(1) − ϕ(r)ϕ(a)) = ϕ(u)ϕ(1 − ra) = ϕ(1) = 1

So ϕ(a) ∈ J(T ) as required.

Theorem:

Let R, T be rings. Then J(R × T ) ⊆ J(R) × J(T ).

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L et π1 : R ×T → R be a (surjective) projection. By the previous theorem, π1 (J(R ×T )) ⊆

J(R).

Similarly let π2 : R × T → T be a (surjective) projection. We see that π2 (J(R × T )) ⊆ J(T ).

**Now take (a, b) ∈ J(R × T ). Note that a = π1 (a, b) ∈ J(R) and b = π2 (a, b) ∈ J(T ). Hence
**

(a, b) ∈ J(R) × J(T ) as required.

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254.6 quasi-regularity

**An element x of a ring is called right quasi-regular [resp. left quasi-regular] if there is
**

an element y in the ring such that x + y + xy = 0 [resp. x + y + yx = 0].

**For calculations with quasi-regularity, it is useful to introduce the operation ∗ defined:
**

x ∗ y = x + y + xy.

Thus x is right quasi-regular if there is an element y such that x ∗ y = 0. The operation ∗ is

easily demonstrated to be associative, and x ∗ 0 = 0 ∗ x = 0 for all x.

**An element x is called quasi-regular if it is both left and right quasi-regular. In this case,
**

there are elements x and y such that x+y+xy = 0 = x+z+zx (equivalently, x∗y = z∗x = 0).

A calculation shows that

y = 0 ∗ y = (z ∗ x) ∗ y = z ∗ (x ∗ y) = z.

So y = z is a unique element, depending on x, called the quasi-inverse of x.

**An ideal (one- or two-sided) of a ring is called quasi-regular if each of its elements is quasi-
**

regular. Similarly, a ring is called quasi-regular if each of its elements is quasi-regular (such

rings cannot have an identity element).

Lemma 1. Let A be an ideal (one- or two-sided) in a ring R. If each element of A is right

quasi-regular, then A is a quasi-regular ideal.

**This lemma means that there is no extra generality gained in defining terms such as right
**

quasi-regular left ideal, etc.

**Quasi-regularity is important because it provides elementary characterizations of the Jacobson radical
**

for rings without an identity element:

• The Jacobson radical of a ring is the sum of all quasi-regular left (or right) ideals.

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• The Jacobson radical of a ring is the largest quasi-regular ideal of the ring.

**For rings with an identity element, note that x is [right, left] quasi-regular if and only if 1 + x
**

is [right, left] invertible in the ring.

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254.7 semiprimitive ring

**A ring R is said to be semiprimitive (sometimes semisimple) if its Jacobson radical is the
**

zero ideal.

Any simple ring is automatically semiprimitive.

**A finite direct product of matrix rings over division rings can be shown to be semiprimitive
**

and both left and right artinian.

**The Artin-Wedderburn Theorem states that any semiprimitive ring which is left or right
**

Artinian is isomorphic to a finite direct product of matrix rings over division rings.

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1120

Chapter 255

**16N40 – Nil and nilpotent radicals,
**

sets, ideals, rings

255.1 Koethe conjecture

The Koethe Conjecture is the statement that for any pair of nil right ideals A and B in any

ring R, the sum A + B is also nil.

**If either of A or B is a two-sided ideal, it is easy to see that A + B is nil. Suppose A is a
**

two-sided ideal, and let x ∈ A + B. The quotient (A + B)/A is nil since it is a homomorphic

image of B. So there is an n > 0 with xn ∈ A. Then there is an m > 0 such that xnm = 0,

because A is nil.

In particular, this means that the Koethe conjecture is true for commutative rings.

It has been shown to be true for many classes of rings, but the general statement is still

unproven, and no counter example has been found.

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255.2 nil and nilpotent ideals

An element x of a ring is nilpotent if xn = 0 for some positive integer n.

**A ring R is nil if every element in R is nilpotent. Similarly, a one- or two-sided ideal is
**

called nil if each of its elements is nilpotent.

**A ring R [resp. a one- or two sided ideal A] is nilpotent if Rn = 0 [resp. An = 0] for some
**

positive integer n.

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A ring or an ideal is locally nilpotent if every finitely generated subring is nilpotent.

The following implications hold for rings (or ideals):

nilpotent ⇒ locally nilpotent ⇒ nil

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1122

Chapter 256

16N60 – Prime and semiprime rings

256.1 prime ring

A ring R is said to be a prime ring if the zero ideal is a prime ideal.

If R is commutative, this is equivalent to being an integral domain.

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1123

Chapter 257

16N80 – General radicals and rings

257.1 prime radical

The prime radical of a ring R is the intersection of all the prime ideals of R.

**Note that the prime radical is the smallest semiprime ideal of R, and that R is a semiprime ring
**

if and only if its prime radical is the zero ideal.

Version: 2 Owner: antizeus Author(s): antizeus

257.2 radical theory

Let x◦ be a property which defines a class of rings, which we will call the x◦ -rings.

Then x◦ is a radical property if it satisfies:

**1. The class of x◦ -rings is closed under homomorphic images.
**

2. Every ring R has a largest ideal in the class of x◦ -rings; this ideal is written x◦ (R).

3. x◦ (R/x◦ (R)) = 0.

**Note: it is extremely important when interpreting the above definition that your definition
**

of a ring does not require an identity element.

**The ideal x◦ (R) is called the x◦ -radical of R. A ring is called x◦ -radical if x◦ (R) = R, and
**

is called x◦ -semisimple if x◦ (R) = 0.

**If x◦ is a radical property, then the class of x◦ -rings is also called the class of x◦ -radical
**

rings.

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The class of x◦ -radical rings is closed under ideal extensions. That is, if A is an ideal of R,

and A and R/A are x◦ -radical, then so is R.

**Radical theory is the study of radical properties and their interrelations. There are several
**

well-known radicals which are of independent interest in ring theory (See examples – to

follow).

**The class of all radicals is however very large. Indeed, it is possible to show that any
**

partition of the class of simple rings into two classes, R and S gives rise to a radical x◦ with

the property that all rings in R are x◦ -radical and all rings in S are x◦ -semisimple.

A radical x◦ is hereditary if every ideal of an x◦ -radical ring is also x◦ -radical.

A radical x◦ is supernilpotent if the class of x◦ -rings contains all nilpotent rings.

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1125

Chapter 258

**16P40 – Noetherian rings and
**

modules

258.1 Noetherian ring

**A ring R is right noetherian (or left noetherian ) if R is noetherian as a right module (or
**

left module ), i.e., if the three equivalent conditions hold:

1. right ideals (or left ideals) are finitely generated

2. the ascending chain condition holds on right ideals (or left ideals)

3. every nonempty family of right ideals (or left ideals) has a maximal element.

**We say that R is noetherian if it is both left noetherian and right noetherian. Examples of
**

Noetherian rings include any field (as the only ideals are 0 and the whole ring) and the ring

Z of integers (each ideal is generated by a single integer, the greatest common divisor of the

elements of the ideal). The Hilbert basis theorem says that a ring R is noetherian iff the

polynomial ring R[x] is.

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258.2 noetherian

A module M is noetherian if it satisfies the following equivalent conditions:

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• the ascending chain condition holds for submodules of M ;

• every nonempty family of submodules of M has a maximal element;

• every submodule of M is finitely generated.

**A ring R is left noetherian if it is noetherian as a left module over itself (i.e. if R R is a
**

noetherian module), and right noetherian if it is noetherian as a right module over itself (i.e.

if RR is an noetherian module), and simply noetherian if both conditions hold.

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1127

Chapter 259

**16P60 – Chain conditions on
**

annihilators and summands:

Goldie-type conditions , Krull

dimension

259.1 Goldie ring

**Let R be a ring. If the set of annihilators {r. ann(x) | x ∈ R} satisifies the ascending chain condition,
**

then R is said to satisfy the ascending chain condition on right annihilators.

**A ring R is called a right Goldie ring if it satisfies the ascending chain condition on right
**

annihilators and RR is a module of finite rank.

**Left Goldie ring is defined similarly. If the context makes it clear on which side the ring
**

operates, then such a ring is simply called a Goldie ring.

A right noetherian ring is right Goldie.

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259.2 uniform dimension

**Let M be a module over a ring R, and suppose that M contains no infinite direct sums of
**

non-zero submodules. (This is the same as saying that M is a module of finite rank.)

1128

Then there exits an integer n such that M contains an essential submodule N where

N = U1 ⊕ U2 ⊕ · · · ⊕ Un

is a direct sum of n uniform submodules.

**This number n does not depend on the choice of N or the decomposition into uniform
**

submodules.

We call n the uniform dimension of M. Sometimes this is written u-dim M = n.

If R is a field K, and M is a finite-dimensional vector space over K, then u-dim M = dimK M.

u-dim M = 0 if and only if M = 0.

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1129

Chapter 260

**16S10 – Rings determined by
**

universal properties (free algebras,

coproducts, adjunction of inverses,

etc.)

260.1 Ore domain

**Let R be a domain. We say that R is a right Ore domain if any two nonzero elements of R
**

have a nonzero common right multiple, i.e. for every pair of nonzero x and y, there exists a

pair of elements r and s of R such that xr = ys 6= 0.

**This condition turns out to be equivalent to the following conditions on R when viewed as
**

a right R-module:

(a) RR is a uniform module.

(b) RR is a module of finite rank.

The definition of a left Ore domain is similar.

If R is a commutative domain, then it is a right (and left) Ore domain.

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Chapter 261

**16S34 – Group rings , Laurent
**

polynomial rings

261.1 support

**Let R[G] be the group ring of a group G over a ring R.
**

P

Let x = g xg g be an element of R[G]. The support of x, often written supp(x), is the set

of elements of G which occur with non-zero coefficient in the expansion of x.

Thus:

supp(x) = {g ∈ G | xg 6= 0}.

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1131

Chapter 262

**16S36 – Ordinary and skew
**

polynomial rings and semigroup rings

262.1 Gaussian polynomials

For an indeterminate u and integers n ≥ m ≥ 0 we define the following:

(a) (m)u = um−1 + um−2 + · · · + 1 for m > 0,

(b) (m!)u = (m)u (m − 1)u · · · (1)u for m > 0, and (0!)u = 1,

n

(n!)u n

(c) m u

= (m!)u ((n−m)!)u

. If m > n then we define m u

= 0.

n

The expressions m u

are called u-binomial coefficients or Gaussian polynomials.

**Note: if we replace u with 1, then we obtain the familiar integers, factorials, and binomial coefficients.
**

Specifically,

(a) (m)1 = m,

(b) (m!)1 = m!,

n

n

(c) m 1

= m .

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262.2 q skew derivation

**Let (σ, δ) be a skew derivation on a ring R. Let q be a central (σ, δ)-constant. Suppose
**

further that δσ = q · σδ. Then we say that (σ, δ) is a q-skew derivation.

Version: 5 Owner: antizeus Author(s): antizeus

262.3 q skew polynomial ring

**If (σ, δ) is a q-skew derivation on R, then we say that the skew polynomial ring R[θ; σ, δ] is
**

a q-skew polynomial ring.

Version: 3 Owner: antizeus Author(s): antizeus

262.4 sigma derivation

**If σ is a ring endomorphism on a ring R, then a (left) σ-derivation is an additive map δ on
**

R such that δ(x · y) = σ(x) · δ(y) + δ(x) · y for all x, y in R.

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262.5 sigma, delta constant

**If (σ, δ) is a skew derivation on a ring R, then a (σ, δ)-constant is an element q of R such
**

that σ(q) = q and δ(q) = 0.

**Note: If q is a (σ, δ)-constant, then it follows that σ(q · x) = q · σ(x) and δ(q · x) = q · δ(x)
**

for all x in R.

Version: 3 Owner: antizeus Author(s): antizeus

262.6 skew derivation

**A (left) skew derivation on a ring R is a pair (σ, δ), where σ is a ring endomorphism of R,
**

and δ is a left σ-derivation on R.

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262.7 skew polynomial ring

**If (σ, δ) is a left skew derivation on R, then we can construct the (left) skew polynomial ring
**

R[θ; σ, δ], which is made up of polynomials in an indeterminate θ and left-hand coefficients

from R, with multiplication satisfying the relation

θ · r = σ(r) · θ + δ(r)

for all r in R.

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1134

Chapter 263

16S99 – Miscellaneous

263.1 algebra

**Let A be a ring with identity. An algebra over A is a ring B with identity together with a
**

ring homomorphism f : A −→ Z(B), where Z(B) denotes the center of B.

Equivalently, an algebra over A is an A–module B which is a ring and satisfies the property

a · (x ∗ y) = (a · x) ∗ y = x ∗ (a · y)

**for all a ∈ A and all x, y ∈ B. Here · denotes A–module multiplication and ∗ denotes
**

ring multiplication in B. One passes between the two definitions as follows: given any ring

homomorphism f : A −→ Z(B), the scalar multiplication rule

a · b := f (a) ∗ b

makes B into an A–module in the sense of the second definition.

Version: 5 Owner: djao Author(s): djao

263.2 algebra (module)

**Given a commutative ring R, an algebra over R is a module M over R, endowed with a law
**

of composition

f :M ×M →M

which is R-bilinear.

**Most of the important algebras in mathematics belong to one or the other of two classes:
**

the unital associative algebras, and the Lie algebras.

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263.2.1 Unital associative algebras

**In these cases, the ”product” (as it is called) of two elements v and w of the module, is
**

denoted simply by vw or v w or the like.

**Any unital associative algebra is an algebra in the sense of djao (a sense which is also used
**

by Lang in his book Algebra (Springer-Verlag)).

Examples of unital associative algebras:

– tensor algebras and – quotients of them

– Cayley algebras, such as the – ring of – quaternions – polynomial rings

**– the ring of – endomorphisms of a – vector space, in which – the bilinear product of – two
**

mappings is simply the – composite mapping.

263.2.2 Lie algebras

In these cases the bilinear product is denoted by [v, w], and satisfies

[v, v] = 0 for all v ∈ M

[v, [w, x]] + [w, [x, v]] + [x, [v, w]] = 0 for all v, w, x ∈ M

The second of these formulas is called the Jacobi identity. One proves easily

[v, w] + [w, v] = 0 for all v, w ∈ M

for any Lie algebra M.

Lie algebras arise naturally from Lie groups, q.v.

Version: 1 Owner: karthik Author(s): Larry Hammick

1136

Chapter 264

16U10 – Integral domains

264.1 Prüfer domain

**An integral domain R is a Prüfer domain if every finitely generated ideal I of R is invert-
**

ible.

Let RI denote the localization of R at I. Then the following statements are equivalent:

• i) R is a Prüfer domain.

• ii) For every prime ideal P in R, RP is a valuation domain.

• iii) For every maximal ideal M in R, RM is a valuation domain.

A Prüfer domain is a Dedekind domain if and only if it is noetherian.

**If R is a Prüfer domain with quotient field K, then any domain S such that R ⊂ S ⊂ K is
**

Prüfer.

REFERENCES

1. Thomas W. Hungerford. Algebra. Springer-Verlag, 1974. New York, NY.

Version: 2 Owner: mathcam Author(s): mathcam

264.2 valuation domain

An integral domain R is a valuation domain if for all a, b ∈ R, either a|b or b|a.

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1138

Chapter 265

**16U20 – Ore rings, multiplicative sets,
**

Ore localization

265.1 Goldie’s Theorem

**Let R be a ring with an identity. Then R has a right classical ring of quotients Q which
**

is semisimple Artinian if and only if R is a semiprime right Goldie ring. If this is the case,

then the composition length of Q is equal to the uniform dimension of R.

**An immediate corollary of this is that a semiprime right noetherian ring always has a right
**

classical ring of quotients.

This result was discovered by Alfred Goldie in the late 1950’s.

Version: 3 Owner: mclase Author(s): mclase

265.2 Ore condition

**A ring R satisfies the left Ore condition (resp. right Ore condition) if and only if for
**

all elements x and y with x regular, there exist elements u and v with v regular such that

ux = vy (resp.xu = yv).

A ring which satisfies the (left, right) Ore condition is called a (left, right) Ore ring.

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265.3 Ore’s theorem

**A ring has a (left, right) classical ring of quotients if and only if it satisfies the (left, right)
**

Ore condition.

Version: 3 Owner: mclase Author(s): mclase

265.4 classical ring of quotients

**Let R be a ring. An element of R is called regular if it is not a right zero divisor or a
**

left zero divisor in R.

**A ring Q ⊃ R is a left classical ring of quotients for R (resp. right classical ring of
**

quotients for R) if it satisifies:

• every regular element of R is invertible in Q

**• every element of Q can be written in the form x−1 y (resp. yx−1 ) with x, y ∈ R and x
**

regular.

If a ring R has a left or right classical ring of quotients, then it is unique up to isomorphism.

**If R is a commutative integral domain, then the left and right classical rings of quotients
**

always exist – they are the field of fractions of R.

For non-commutative rings, necessary and sufficient conditions are given by Ore’s theorem.

**Note that the goal here is to construct a ring which is not too different from R, but in
**

which more elements are invertible. The first condition says which elements we want to be

invertible. The second condition says that Q should contain just enough extra elements to

make the regular elements invertible.

**Such rings are called classical rings of quotients, because there are other rings of quotients.
**

These all attempt to enlarge R somehow to make more elements invertible (or sometimes to

make ideals invertible).

Finally, note that a ring of quotients is not the same as a quotient ring.

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265.5 saturated

Let S be multiplicative subset of A. We say that S is a saturated if

ab ∈ S ⇒ a, b ∈ S.

**When A is an integral domain, then S is saturated if and only if its complement A\S is
**

union of prime ideals.

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1141

Chapter 266

**16U70 – Center, normalizer (invariant
**

elements)

266.1 center (rings)

**If A is a ring, the center of A, sometimes denoted Z(A), is the set of all elements in A that
**

commute with all other elements of A. That is,

Z(A) = {a ∈ A | ax = xa∀x ∈ A}

**Note that 0 ∈ Z(A) so the center is non-empty. If we assume that A is a ring with a
**

multiplicative unity 1, then 1 is in the center as well. The center of A is also a subring of A.

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Chapter 267

16U99 – Miscellaneous

267.1 anti-idempotent

**An element x of a ring is called an anti-idempotent element, or simply an anti-idempotent
**

if x2 = −x.

**The term is most often used in linear algebra. Every anti-idempotent matrix over a field is
**

diagonalizable. Two anti-idempotent matrices are similar if and only if they have the same

rank.

Version: 1 Owner: mathcam Author(s): mathcam

1143

Chapter 268

**16W20 – Automorphisms and
**

endomorphisms

268.1 ring of endomorphisms

Let R be a ring and let M be a right R-module.

**An endomorphism of M is a R-module homomorphism from M to itself. We shall write
**

endomorphisms on the left, so that f : M → M maps x 7→ f (x). If f, g : M → M are two

endomorphisms, we can add them:

f + g : x 7→ f (x) + g(x)

and multiply them

f g : x 7→ f (g(x))

With these operations, the set of endomorphisms of M becomes a ring, which we call the

ring of endomorphisms of M, written EndR (M).

**Instead of writing endomorphisms as functions, it is often convenient to write them multi-
**

plicatively: we simply write the application of the endomorphism f as x 7→ f x. Then the

fact that each f is an R-module homomorphism can be expressed as:

f (xr) = (f x)r

for all x ∈ M and r ∈ R and f ∈ EndR (M). With this notation, it is clear that M becomes

an EndR (M)-R-bimodule.

**Now, let N be a left R-module. We can construct the ring EndR (N) in the same way. There
**

is a complication, however, if we still think of endomorphism as functions written on the

left. In order to make M into a bimodule, we need to define an action of EndR (N) on the

right of N: say

x · f = f (x)

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But then we have a problem with the multiplication:

x · f g = f g(x) = f (g(x))

but

(x · f ) · g = f (x) · g = g(f (x))!

In order to make this work, we need to reverse the order of composition when we define

multiplication in the ring EndR (N) when it acts on the right.

**There are essentially two different ways to go from here. One is to define the muliplication in
**

EndR (N) the other way, which is most natural if we write the endomorphisms as functions

on the right. This is the approach taken in many older books.

**The other is to leave the muliplication in EndR (N) the way it is, but to use the opposite ring
**

to define the bimodule. This is the approach that is generally taken in more recent works.

Using this approach, we conclude that N is a R-EndR (N)op -bimodule. We will adopt this

convention for the lemma below.

**Considering R as a right and a left module over itself, we can construct the two endomor-
**

phism rings EndR (RR ) and EndR (R R).

**Lemma 2. Let R be a ring with an identity element. Then R ' EndR (RR ) and R '
**

EndR (R R)op .

D efine ρr ∈ EndR (R R) by x 7→ xr.

**A calculation shows that ρrs = ρs ρr (functions written on the left) from which it is easily
**

seen that the map θ : r 7→ ρr is a ring homomorphism from R to EndR (R R)op .

We must show that this is an isomorphism.

If ρr = 0, then r = 1r = ρr (1) = 0. So θ is injective.

**Let f be an arbitrary element of EndR (R R), and let r = f (1). Then for any x ∈ R,
**

f (x) = f (x1) = xf (1) = xr = ρr (x), so f = ρr = θ(r).

The proof of the other isomorphism is similar.

Version: 4 Owner: mclase Author(s): mclase

1145

Chapter 269

**16W30 – Coalgebras, bialgebras, Hopf
**

algebras ; rings, modules, etc. on

which these act

269.1 Hopf algebra

A Hopf algebra is a bialgebra A over a field K with a K-linear map S : A → A, called the

Definition 1. antipode, such that

m ◦ (S ⊗ id) ◦ ∆ = η ◦ ε = m ◦ (id ⊗ S) ◦ ∆, (269.1.1)

**where m : A ⊗ A → A is the multiplication map m(a ⊗ b) = ab and η : K → A is the unit
**

map η(k) = k1I.

In terms of a commutative diagram:

A

∆ ∆

A⊗A ε A⊗A

S⊗id C id⊗∗S

A⊗A η A⊗A

m m

A

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Example 1 (Algebra of functions on a finite group). Let A = C(G) be the algebra of complex-

valued functions on a finite group G and identify C(G × G) with A ⊗ A. Then, A is a

Hopf algebra with comultiplication (∆(f ))(x, y) = f (xy), counit ε(f ) = f (e), and antipode

(S(f ))(x) = f (x−1 ).

Example 2 (Group algebra of a finite group). Let A = CG be the complex group algebra

of a finite group G. Then, A is a Hopf algebra with comultiplication ∆(g) = g ⊗ g, counit

ε(g) = 1, and antipode S(g) = g −1 .

**The above two examples are dual to one another. Define a bilinear form C(G) ⊗ CG → C
**

by hf, xi = f (x). Then,

hf g, xi = hf ⊗ g, ∆(x)i,

h1, xi = ε(x),

h∆(f ), x ⊗ yi = hf, xyi,

ε(f ) = hf, ei,

hS(f ), xi = hf, S(x)i.

Example 3 (Polynomial functions on a Lie group). Let A = Poly(G) be the algebra of

complex-valued polynomial functions on a complex Lie group G and identify Poly(G × G)

with A ⊗ A. Then, A is a Hopf algebra with comultiplication (∆(f ))(x, y) = f (xy), counit

ε(f ) = f (e), and antipode (S(f ))(x) = f (x−1 ).

Example 4 (Universal enveloping algebra of a Lie algebra). Let A = U(g) be the universal enveloping algebra

of a complex Lie algebra g. Then, A is a Hopf algebra with comultiplication ∆(X) =

X ⊗ 1 + 1 ⊗ X, counit ε(X) = 0, and antipode S(X) = −X.

**The above two examples are dual to one another (if g is the Lie algebra of G). Define a
**

d

bilinear form Poly(G) ⊗ U(g) → C by hf, Xi = dt t=0

f (exp(tX)).

Version: 6 Owner: mhale Author(s): mhale

269.2 almost cocommutative bialgebra

**A bialgebra A is called almost cocommutative if there is an unit R ∈ A ⊗ A such that
**

R∆(a) = ∆op (a)R

where ∆op is the opposite comultiplication (the usual comultiplication, composed with the

flip map of the tensor product A ⊗ A). The element R is often called the R-matrix of A.

**The significance of the almost cocommutative condition is that σV,W = σ ◦ R : V ⊗ W →
**

W ⊗ V gives a natural isomorphism of bialgebra representations, where V and W are A-

modules, making the category of A-modules into a quasi-tensor or braided monoidal category.

Note that σW,V ◦ σV,W is not necessarily the identity (this is the braiding of the category).

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269.3 bialgebra

A

Definition 2. bialgebra is a vector space that is both a unital algebra and a coalgebra, such

that the comultiplication and counit are unital algebra homomorphisms.

Version: 2 Owner: mhale Author(s): mhale

269.4 coalgebra

A

Definition 3. coalgebra is a vector space A over a field K with a K-linear map ∆ : A →

A ⊗ A, called the

Definition 4. comultiplication, and a (non-zero) K-linear map ε : A → K, called the

Definition 5. counit, such that

(∆ ⊗ id) ◦ ∆ = (id ⊗ ∆) ◦ ∆ (coassociativity), (269.4.1)

(ε ⊗ id) ◦ ∆ = id = (id ⊗ ε) ◦ ∆. (269.4.2)

**In terms of commutative diagrams:
**

A

∆ ∆

A⊗A A⊗A

∆⊗id id⊗∆

A⊗A⊗A

A

∆ ∆

A⊗A id A⊗A

ε⊗id id⊗ε

A

**Let σ : A ⊗ A → A ⊗ A be the flip map σ(a ⊗ b) = b ⊗ a. A coalgebra is said to be
**

Definition 6. cocommutative if σ ◦ ∆ = ∆.

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269.5 coinvariant

**Let V be a comodule with a right coaction t : V → V ⊗ A of a coalgebra A. An element
**

v ∈ V is

Definition 7. right coinvariant if

t(v) = v ⊗ 1IA . (269.5.1)

Version: 1 Owner: mhale Author(s): mhale

269.6 comodule

Let (A, ∆, ε) be a coalgebra. A

Definition 8. right A-comodule is a vector space V with a linear map t : V → V ⊗A, called

the

Definition 9. right coaction, satisfying

(t ⊗ id) ◦ t = (id ⊗ ∆) ◦ t, (id ⊗ ε) ◦ t = id. (269.6.1)

An A-comodule is also referred to as a corepresentation of A.

**Let V and W be two right A-comodules. Then V ⊕ W is also a right A-comodule. If A is
**

a bialgebra then V ⊗ W is a right A-comodule as well (make use of the multiplication map

A ⊗ A → A).

Version: 2 Owner: mhale Author(s): mhale

269.7 comodule algebra

**Let H be a bialgebra. A right H-comodule algebra is a unital algebra A which is a right
**

H-comodule satisfying

X

t(ab) = t(a)t(b) = a(1) b(1) ⊗ a(2) b(2) , t(1IA ) = 1IA ⊗ 1IH , (269.7.1)

for all h ∈ H and a, b ∈ A.

**There is a dual notion of a H-module coalgebra.
**

Example 5. Let H be a bialgebra. Then H is itself a H-comodule algebra for the right

regular coaction t(h) = ∆(h).

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269.8 comodule coalgebra

**Let H be a bialgebra. A right H-comodule coalgebra is a coalgebra A which is a right
**

H-comodule satisfying

X

(∆ ⊗ id)t(a) = a(1)(1) ⊗ a(2)(1) ⊗ a(1)(2) a(2)(2) , (ε ⊗ id)t(a) = ε(a)1IH , (269.8.1)

for all h ∈ H and a ∈ A.

**There is a dual notion of a H-module algebra.
**

Example 6. Let H be a Hopf algebra. Then H is itself a H-comodule coalgebra for the

adjoint coaction t(h) = h(2) ⊗ S(h(1) )h(3) .

Version: 4 Owner: mhale Author(s): mhale

269.9 module algebra

**Let H be a bialgebra. A left H-module algebra is a unital algebra A which is a left
**

H-module satisfying

X

h . (ab) = (h(1) . a)(h(2) . b), h . 1IA = ε(h)1IA , (269.9.1)

for all h ∈ H and a, b ∈ A.

**There is a dual notion of a H-comodule coalgebra.
**

ExampleP7. Let H be a Hopf algebra. Then H is itself a H-module algebra for the adjoint action

g . h = g(1) hS(g(2) ).

Version: 4 Owner: mhale Author(s): mhale

269.10 module coalgebra

**Let H be a bialgebra. A left H-module coalgebra is a coalgebra A which is a left H-
**

module satisfying

X

∆(h . a) = (h(1) . a(1) ) ⊗ (h(2) . a(2) ), ε(h . a) = ε(h)ε(a), (269.10.1)

for all h ∈ H and a ∈ A.

**There is a dual notion of a H-comodule algebra.
**

Example 8. Let H be a bialgebra. Then H is itself a H-module coalgebra for the left regular

action g . h = gh.

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1150

Chapter 270

16W50 – Graded rings and modules

270.1 graded algebra

An algebra A is graded if it is a graded module and satisfies

Ap · Aq ⊆ Ap+q

**Examples of graded algebras include the polynomial ring k[X] being an N-graded k-algebra,
**

and the exterior algebra.

Version: 1 Owner: dublisk Author(s): dublisk

270.2 graded module

**If R = R0 ⊕ R1 ⊕ · · · is a graded ring, then a graded module over R is a module M of the
**

form M = ⊕∞ i=−∞ Mi and satisfies Ri Mj ⊆ Mi+j for all i, j.

Version: 4 Owner: KimJ Author(s): KimJ

270.3 supercommutative

**Let R be a Z2 -graded ring. Then R is supercommutative if for any homogeneous elements a
**

and b ∈ R:

ab = (−1)deg a deg b ba.

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This is, even homogeneous elements are in the center of the ring, and odd homogeneous

elements anti-commute.

**Common examples of supercommutative rings are the exterior algebra of a module over a
**

commutative ring (in particular, a vector space) and the cohomology ring of a topological space

(both with the standard grading by degree reduced mod 2).

Version: 1 Owner: bwebste Author(s): bwebste

1152

Chapter 271

**16W55 – “Super” (or “skew”)
**

structure

271.1 super tensor product

**If A and B are Z-graded algebras, we define the super tensor product A ⊗su B to be the
**

ordinary tensor product as graded modules, but with multiplication - called the super product

- defined by

0

(a ⊗ b)(a0 ⊗ b0 ) = (−1)(deg b)(deg a ) aa0 ⊗ bb0

where a, a0 , b, b0 are homogeneous. The super tensor product of A and B is itself a graded algebra,

as we grade the super tensor product of A and B as follows:

a

(A ⊗su B)n = Ap ⊗ B q

p,q : p+q=n

Version: 4 Owner: dublisk Author(s): dublisk

271.2 superalgebra

A graded algebra A is said to be a super algebra if it has a Z/2Z grading.

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271.3 supernumber

**Let ΛN be the Grassmann algebra generated by θi , i = 1 . . . N, such that θi θj = −θj θi and
**

(θi )2 = 0. Denote by Λ∞ , the case of an infinite number of generators θi . A

Definition 10. supernumber is an element of ΛN or Λ∞ .

**Any supernumber z can be expressed uniquely in the form
**

1 1

z = z0 + zi θi + zij θi θj + . . . + zi1 ...in θi1 . . . θin + . . . ,

2 n!

where the coefficients zi1 ...in ∈ C are antisymmetric in their indices. The

Definition 11. body of z is defined as zB = z0 , and its

**Definition 12. soul is defined as zS = z − zB . If zB 6= 0 then z has an inverse given by
**

k

−1 1 X zS

z = − .

zB k=0 zB

**A supernumber can be decomposed into the even and odd parts
**

1 1

zeven = z0 + zij θi θj + . . . + zi1 ...i2n θi1 . . . θi2n + . . . ,

2 (2n)!

1 1

zodd = zi θi + zijk θi θj θk + . . . + zi ...i θi1 . . . θi2n+1 + . . . .

6 (2n + 1)! 1 2n+1

Purely even supernumbers are called

Definition 13. c-numbers, and odd supernumbers are called

**Definition 14. a-numbers. The superalgebra ΛN thus has a decomposition ΛN = Cc ⊕ Ca ,
**

where Cc is the space of c-numbers, and Ca is the space of a-numbers.

**Supernumbers are the generalisation of complex numbers to a commutative superalgebra of
**

commuting and anticommuting “numbers”. They are primarily used in the description of

fermionic fields in quantum field theory.

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1154

Chapter 272

16W99 – Miscellaneous

272.1 Hamiltonian quaternions

Definition of Q

**We define a unital associative algebra Q over R, of dimension 4, by the basis {1, i, j, k} and
**

the multiplication table

1 i j k

i −1 k −j

j −k −1 i

k j −i −1

**(where the element in row x and column y is xy, not yx). Thus an arbitrary element of Q
**

is of the form

a1 + bi + cj + dk, a, b, c, d ∈ R

(sometimes denoted by ha, b, c, di or by a+hb, c, di) and the product of two elements ha, b, c, di

and hα, β, γ, δi is hw, x, y, zi where

w = aα − bβ − cγ − dδ

x = aβ + bα + cδ − dγ

y = aγ − bδ + cα + kβ

z = aδ + bγ − cβ + kα

The elements of Q are known as Hamiltonian quaternions.

**Clearly the subspaces of Q generated by {1} and by {1, i} are subalgebras isomorphic to R
**

and C respectively. R is customarily identified with the corresponding subalgebra of Q. (We

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shall see in a moment that there are other and less obvious embeddings of C in Q.) The

real numbers commute with all the elements of Q, and we have

λ · ha, b, c, di = hλa, λb, λc, λdi

for λ ∈ R and ha, b, c, di ∈ Q.

norm, conjugate, and inverse of a quaternion

**Like the complex numbers (C), the quaternions have a natural involution called the quater-
**

nion conjugate. If q = a1 + bi + cj + dk, then the quaternion conjugate of q, denoted q, is

simply q = a1 − bi − cj − dk.

**One can readily verify that if q = a1 + bi + cj + dk, then qq = (a2 + b2 + c2 + d2 )1. (See
**

√ This product is used to form a norm | · | on the algebra (or the

Euler four-square identity.)

ring) Q: We define kqk = s where qq = s1.

If v, w ∈ Q and λ ∈ R, then

1. kvk > 0 with equality only if v = h0, 0, 0, 0i = 0

2. kλvk = |λ|kvk

3. kv + wk 6 kvk + kwk

4. kv · wk = kvk · kwk

which means that Q qualifies as a normed algebra when we give it the norm | · |.

**Because the norm of any nonzero quaternion q is real and nonzero, we have
**

qq qq

2

= = h1, 0, 0, 0i

kqk kqk2

which shows that any nonzero quaternion has an inverse:

q

q −1 = .

kqk2

Other embeddings of C into Q

**One can use any non-zero q to define an embedding of C into Q. If n(z) is a natural embedding
**

of z ∈ C into Q, then the embedding:

z → qn(z)q −1

is also an embedding into Q. Because Q is an associative algebra, it is obvious that:

(qn(a)q −1 )(qn(b)q −1 ) = q(n(a)n(b))q −1

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and with the distributive laws, it is easy to check that

(qn(a)q −1 ) + (qn(b)q −1 ) = q(n(a) + n(b))q −1

Rotations in 3-space

Let us write

U = {q ∈ Q : ||q|| = 1}

With multiplication, U is a group. Let us briefly sketch the relation between U and the

group SO(3) of rotations (about the origin) in 3-space.

**An arbitrary element q of U can be expressed cos 2θ + sin θ2 (ai + bj + ck), for some real
**

numbers θ, a, b, c such that a2 + b2 + c2 = 1. The permutation v 7→ qv of U thus gives rise to

a permutation of the real sphere. It turns out that that permutation is a rotation. Its axis

is the line through (0, 0, 0) and (a, b, c), and the angle through which it rotates the sphere

is θ. If rotations F and G correspond to quaternions q and r respectively, then clearly the

permutation v 7→ qrv corresponds to the composite rotation F ◦ G. Thus this mapping of

U onto SO(3) is a group homomorphism. Its kernel is the subset {1, −1} of U, and thus it

comprises a double cover of SO(3). The kernel has a geometric interpretation as well: two

unit vectors in opposite directions determine the same axis of rotation.

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1157

Chapter 273

16Y30 – Near-rings

273.1 near-ring

**A near-ring is a set N together with two binary operations, denoted + : N × N → N and
**

· : N × N → N, such that

**1. (a + b) + c = a + (b + c) and (a · b) · c = a · (b · c) for all a, b, c ∈ N (associativity of
**

both operations)

**2. There exists an element 0 ∈ N such that a + 0 = 0 + a = a for all a ∈ N (additive
**

identity)

3. For all a ∈ N, there exists b ∈ N such that a + b = b + a = 0 (additive inverse)

4. (a + b) · c = (a · c) + (b · c) for all a, b, c ∈ N (right distributive law)

Note that the axioms of a near-ring differ from those of a ring in that they do not require

addition to be commutative, and only require distributivity on one side.

Every element a in a near-ring has a unique additive inverse, denoted −a.

**We say N has an identity element if there exists an element 1 ∈ N such that a · 1 = 1 · a = a
**

for all a ∈ N. We say N is distributive if a · (b + c) = (a · b) + (a · c) holds for all a, b, c ∈ N.

We say N is commutative if a · b = b · a for all a, b ∈ N.

**A natural example of a near-ring is the following. Let (G, +) be a group (not necessarily
**

abelian), and let M be the set of all functions from G to G. For two functions f and g in M

define f + g ∈ M by (f + g)(x) = f (x) + g(x) for all x ∈ G. Then (M, +, ◦) is a near-ring

with identity, where ◦ denotes composition of functions.

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1158

Chapter 274

17A01 – General theory

274.1 commutator bracket

**Let A be an associative algebra over a field K. For a, b ∈ A, the element of A defined by
**

[a, b] = ab − ba

is called the commutator of a and b. The corresponding bilinear operation

[−, −] : A × A → A

is called the commutator bracket.

**The commutator bracket is bilinear, skew-symmetric, and also satisfies the Jacobi identity.
**

To wit, for a, b, c ∈ A we have

[a, [b, c]] + [b, [c, a]] + [c, [a, b]] = 0.

The proof of this assertion is straightforward. Each of the brackets in the left-hand side

expands to 4 terms, and then everything cancels.

**In categorical terms, what we have here is a functor from the category of associative algebras
**

to the category of Lie algebras over a fixed field. The action of this functor is to turn an

associative algebra A into a Lie algebra that has the same underlying vector space as A, but

whose multiplication operation is given by the commutator bracket. It must be noted that

this functor is right-adjoint to the universal enveloping algebra functor.

Examples

**• Let V be a vector space. Composition endows the vector space of endomorphisms
**

End V with the structure of an associative algebra. However, we could also regard

End V as a Lie algebra relative to the commutator bracket:

[X, Y ] = XY − Y X, X, Y ∈ End V.

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• The algebra of differential operators has some interesting properties when viewed as a

Lie algebra. The fact is that even though, even though the composition of differential

operators is a non-commutative operation, it is commutative when restricted to the

highest order terms of the involved operators. Thus, if X, Y are differential operators

of order p and q, respectively, the compositions XY and Y X have order p + q. Their

highest order term coincides, and hence the commutator [X, Y ] has order p + q − 1.

**• In light of the preceding comments, it is evident that the vector space of first-order
**

differential operators is closed with respect to the commutator bracket. Specializing

even further we remark that, a vector field is just a homogeneous first-order differential

operator, and that the commutator bracket for vector fields, when viewed as first-order

operators, coincides with the usual, geometrically motivated vector field bracket.

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1160

Chapter 275

17B05 – Structure theory

275.1 Killing form

**Let g be a finite dimensional Lie algebra over a field k, and adX : g → g be the adjoint action,
**

adX Y = [X, Y ]. Then the Killing form on g is a bilinear map

Bg : g × g → k

given by

Bg(X, Y ) = tr(adX ◦ adY ).

The Killing form is invariant and symmetric (since trace is symmetric).

Version: 4 Owner: bwebste Author(s): bwebste

275.2 Levi’s theorem

**Let g be a complex Lie algebra, r its radical. Then the extension 0 → r → g → g/r → 0 is
**

split, i.e., there exists a subalgebra h of g mapping isomorphically to g/r under the natural

projection.

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275.3 nilradical

**Let g be a Lie algebra. Then the nilradical n of g is defined to be the intersection
**

T of

the kernels of all the irreducible representations of g. Equivalently, n = [g, g] rad g, the

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interesection of the derived ideal and radical of g.

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275.4 radical

**Let g be a Lie algebra. Since the sum of any two solvable ideals of g is in turn solvable, there
**

is a unique maximal solvable ideal of any Lie algebra. This ideal is called the radical of g.

Note that g/rad g has no solvable ideals, and is thus semi-simple. Thus, every Lie algebra is

an extension of a semi-simple algebra by a solvable one.

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Chapter 276

**17B10 – Representations, algebraic
**

theory (weights)

276.1 Ado’s theorem

**Every finite dimensional Lie algebra has a faithful finite dimensional representation. In other
**

words, every finite dimensional Lie algebra is a matrix algebra.

This result is not true for Lie groups.

Version: 2 Owner: bwebste Author(s): bwebste

276.2 Lie algebra representation

A representation of a Lie algebra g is a Lie algebra homomorphism

ρ : g → End V,

**where End V is the commutator Lie algebra of some vector space V . In other words, ρ is a
**

linear mapping that satisfies

ρ([a, b]) = ρ(a)ρ(b) − ρ(b)ρ(a), a, b ∈ g

Alternatively, one calls V a g-module, and calls ρ(a), a ∈ g the action of a on V .

We call the representation faithful if ρ is injective.

**A invariant subsspace or sub-module W ⊂ V is a subspace of V satisfying ρ(a)(W ) ⊂ W for
**

all a ∈ g. A representation is called irreducible or simple if its only invariant subspaces are

{0} and the whole representation.

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The dimension of V is called the dimension of the representation. If V is infinite-dimensional,

then one speaks of an infinite-dimensional representation.

**Given a representation or pair of representation, there are a couple of operations which will
**

produce other representations:

**First there is direct sum. If ρ : g → End(V ) and σ : g → End(W ) are representations,
**

then V ⊕ W has the obvious Lie algebra action, by the embedding End(V ) × End(W ) ,→

End(V ⊕ W ).

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276.3 adjoint representation

**Let g be a Lie algebra. For every a ∈ g we define the adjoint endomorphism, a.k.a. the
**

adjoint action,

ad(a) : g → g

to be the linear transformation with action

ad(a) : b 7→ [a, b], b ∈ g.

The linear mapping ad : g → End(g) with action

a 7→ ad(a), a∈g

**is called the adjoint representation of g. The fact that ad defines a representation is a
**

straight-forward consequence of the Jacobi identity axiom. Indeed, let a, b ∈ g be given. We

wish to show that

ad([a, b]) = [ad(a), ad(b)],

where the bracket on the left is the g multiplication structure, and the bracket on the right

is the commutator bracket. For all c ∈ g the left hand side maps c to

[[a, b], c],

while the right hand side maps c to

[a, [b, c]] + [b, [a, c]].

**Taking skew-symmetry of the bracket as a given, the equality of these two expressions is
**

logically equivalent to the Jacobi identity:

[a, [b, c]] + [b, [c, a]] + [c, [a, b]] = 0.

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276.4 examples of non-matrix Lie groups

While most well-known Lie groups are matrix groups, there do in fact exist Lie groups which

are not matrix groups. That is, they have no faithful finite dimensional representations.

For example, let H be the real Heisenberg group

1 a b

H = 0 1 c | a, b, c ∈ R ,

0 0 1

**and Γ the discrete subgroup
**

1 0 n

Γ = 0 1 0 | n ∈ Z .

0 0 1

**The subgroup Γ is central, and thus normal. The Lie group H/Γ has no faithful finite
**

dimensional representations over R or C.

**Another example is the universal cover of SL2 R. SL2 R is homotopy equivalent to a circle,
**

and thus π(SL2 R) ∼ = Z, and thus has an infinite-sheeted cover. Any real or complex repre-

sentation of this group factors through the projection map to SL2 R.

Version: 3 Owner: bwebste Author(s): bwebste

276.5 isotropy representation

**Let g be a Lie algebra, and h ⊂ g a subalgebra. The isotropy representation of h relative to
**

g is the naturally defined action of h on the quotient vector space g/h.

Here is a synopsis of the technical details. As is customary, we will use

b + h, b ∈ g

**to denote the coset elements of g/h. Let a ∈ h be given. Since h is invariant with respect to
**

adg(a), the adjoint action factors through the quotient to give a well defined endomorphism

of g/h. The action is given by

b + h 7→ [a, b] + h, b ∈ g.

This is the action alluded to in the first paragraph.

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Chapter 277

**17B15 – Representations, analytic
**

theory

277.1 invariant form (Lie algebras)

**Let V be a representation of a Lie algebra g over a field k. Then a bilinear form B : V ×V →
**

k is invariant if

B(Xv, w) + B(v, Xw) = 0.

for all X ∈ g, v, w ∈ V . This criterion seems a little odd, but in the context of Lie algebras,

it makes sense. For example, the map B̃ : V → V ∗ given by v 7→ B(·, v) is equivariant if and

only if B is an invariant form.

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Chapter 278

**17B20 – Simple, semisimple,
**

reductive (super)algebras (roots)

278.1 Borel subalgebra

**Let g be a semi-simple Lie group, h a Cartan subalgebra, R the associated root system and
**

R+ ⊂ R a set of positive roots. We have a root decomposition into the Cartan subalgebra

and the root spaces gα !

M

g=h⊕ gα .

α∈R

**Now let b be the direct sum of the Cartan subalgebra and the positive root spaces.
**

M

b=h⊕ gβ .

β∈R+

This is called a Borel subalgebra.

Version: 2 Owner: bwebste Author(s): bwebste

278.2 Borel subgroup

**Let G be a complex semi-simple Lie group. Then any maximal solvable subgroup B 6 G
**

is called a Borel subgroup. All Borel subgroups of a given group are conjugate. Any Borel

group is connected and equal to its own normalizer, and contains a unique Cartan subgroup.

The intersection of B with a maximal compact subgroup K of G is the maximal torus of K.

If G = SLn C, then the standard Borel subgroup is the set of upper triangular matrices.

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278.3 Cartan matrix

**Let R ⊂ E be a reduced root system, with E a euclidean vector space, with inner product
**

(·, ·), and let Π = {α1 , · · · , αn } be a base of this root system. Then the Cartan matrix of

the root system is the matrix

2(αi , αj )

Ci,j = .

(αi , αi )

The Cartan matrix uniquely determines the root system, and is unique up to simultaneous

permutation of the rows and columns. It is also the basis change matrix from the basis of

fundamental weights to the basis of simple roots in E.

Version: 1 Owner: bwebste Author(s): bwebste

278.4 Cartan subalgebra

**Let g be a Lie algebra. Then a Cartan subalgebra is a maximal subalgebra of g which is self-
**

normalizing, that is, if [g, h] ∈ h for all h ∈ h, then g ∈ h as well. Any Cartan subalgebra

h is nilpotent, and if g is semi-simple, it is abelian. All Cartan subalgebras of a Lie algebra

are conjugate by the adjoint action of any Lie group with algebra g.

Version: 3 Owner: bwebste Author(s): bwebste

278.5 Cartan’s criterion

A Lie algebra g is semi-simple if and only if its Killing form Bg is nondegenerate.

Version: 2 Owner: bwebste Author(s): bwebste

278.6 Casimir operator

**Let g be a semisimple Lie algebra, and let (·, ·) denote the Killing form. If {gi} is a basis of
**

g, then there is a dual

P basis {g i} with respect to the Killing form, i.e., (gi , g j ) = δij . Consider

the element Ω = gi g i of the universal enveloping algebra of g. This element, called the

Casimir operator is central in the enveloping algebra, and thus commutes with the g action

on any representation.

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278.7 Dynkin diagram

**Dynkin diagrams are a combinatorial way of representing the imformation in a root system.
**

Their primary advantage is that they are easier to write down, remember, and analyze than

explicit representations of a root system. They are an important tool in the classification of

simple Lie algebras.

**Given a reduced root system R ⊂ E, with E an inner-product space, choose a base or
**

simple roots Π (or equivalently, a set of positive roots R+ ). The Dynkin diagram associated

to R is a graph whose vertices are Π. If πi and πj are distinct elements of the root system, we

−4(πi ,πj )2

add mij = (πi ,πi)(π j ,πj )

lines between them. This number is obivously positive, and an integer

since it is the product of 2 quantities that the axioms of a root system require to be integers.

By the Cauchy-Schwartz inequality, and the fact that simple roots are never anti-parallel

(they are all strictly contained in some half space), mij ∈ {0, 1, 2, 3}. Thus Dynkin diagrams

are finite graphs, with single, double or triple edges. Fact, the criteria are much stronger

than this: if the multiple edges are counted as single edges, all Dynkin diagrams are trees,

and have at most one multiple edge. In fact, all Dynkin diagrams fall into 4 infinite families,

and 5 exceptional cases, in exact parallel to the classification of simple Lie algebras.

**(Does anyone have good Dynkin diagram pictures? I’d love to put some up, but am decidedly
**

lacking.)

Version: 1 Owner: bwebste Author(s): bwebste

278.8 Verma module

**Let g be a semi-simple Lie algebra, h a Cartan subalgebra, and b a Borel subalgebra. Let
**

Fλ for a weight λ ∈ h∗ be the 1-d dimensional b module on which h acts by multiplication

by λ, and the positive root spaces act trivially. Now, the Verma module Mλ of the weight λ

is the g module

Mλ = Fλ ⊗U(b) U(g).

**This is an infinite dimensional representation, and it has a very important property: If V
**

is any representation with highest weight λ, there is a surjective homomorphism Mλ → V .

That is, all representations with highest weight λ are quotients of Mλ . Also, Mλ has a unique

maximal submodule, so there is a unique irreducible representation with highest weight λ.

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278.9 Weyl chamber

**If R ⊂ E is a root system, with E a euclidean vector space, and R+ is a set of positive roots,
**

then the positive Weyl chamber is the set

C = {e ∈ E|(e, α) > 0 ∀α ∈ R+ }.

**The interior of C is a fundamental domain for the action of the Weyl group on E. The image
**

w(C) of C under the any element of the Weyl group is called a Weyl chamber. The Weyl

group W acts simply transitively on the set of Weyl chambers.

A weight which lies inside the positive Weyl chamber is called dominant

Version: 2 Owner: bwebste Author(s): bwebste

278.10 Weyl group

**The Weyl group WR of a root system R ⊂ E, where E is a euclidean vector space, is the
**

subgroup of GL(E) generated by reflection in the hyperplanes perpendicular to the roots.

The map of reflection in a root α is given by

(v, α)

rα (v) = v − 2

(α, α)

.

**The Weyl group is generated by reflections in the simple roots for any choice of a set of
**

positive roots. There is a well-defined length function ` : WR → Z, where `(w) is the

minimal number of reflections in simple roots that w can be written as. This is also the

number of positive roots that w takes to negative roots.

Version: 1 Owner: bwebste Author(s): bwebste

278.11 Weyl’s theorem

**Let g be a finite dimensional semi-simple Lie algebra. Then any finite dimensional representation
**

of g is completely reducible.

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278.12 classification of finite-dimensional representa-

tions of semi-simple Lie algebras

**If g is a semi-simple Lie algebra, then we say that an irreducible representation V has highest
**

weight λ, if there is a vector v ∈ Vλ , the weight space of λ, such that Xv = 0 for X in any

positive root space, and v is called a highest vector, or vector of highest weight.

**There is a unique (up to isomorphism) irreducible finite dimensional representation of g with
**

highest weight λ for any dominant weight λ ∈ ΛW , where ΛW is the weight lattice of g, and

every irreducible representation of g is of this type.

Version: 1 Owner: bwebste Author(s): bwebste

278.13 cohomology of semi-simple Lie algebras

There are some important facts that make the cohomology of semi-simple Lie algebras easier

to deal with than general Lie algebra cohomology.

**In particular, there are a number of vanishing theorems. First of all, let g be a finite-dimensional,
**

semi-simple Lie algebra over C.

Theorem. Let M be an irreducible representation of g. Then H n (g, M) = 0 for all n.

Whitehead’s lemmata. Let M be any representation of g, then H 1 (g, M) = H 2 (g, M) = 0.

**Whitehead’s lemmata lead to two very important results. From the vanishing of H 1 , we
**

can derive Weyl’s theorem, the fact that representations of semi-simple Lie algebras are

completely reducible, since extensions of M by N are classified by H 1 (g, HomMN). And

from the vanishing of H 2 , we obtain Levi’s theorem, which states that every Lie algebra

is a split extension of a semi-simple algebra by a solvable algebra since H 2 (g, M) classifies

extensions of g by M with a specified action of g on M.

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278.14 nilpotent cone

**Let g be a finite dimensional semisimple Lie algebra. Then the nilpotent cone N of g is set
**

of elements which act nilpotently on all representations of g. This is a irreducible subvariety

of g (considered as a k-vector space), which is invariant under the adjoint action of G on g

(here G is the adjoint group associated to g).

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278.15 parabolic subgroup

**Let G be a complex semi-simple Lie group. Then any subgroup P of G containg a Borel subgroup
**

B is called parabolic. Parabolics are classified in the following manner. Let g be the

Lie algebra of G, h the unique Cartan subalgebra contained in b, the algebra of B, R the

set of roots corresponding to this choice of Cartan, and R+ the set of positive roots whose

root spaces are contained in b and let p be the Lie algebra of P . Then there exists a unique

subset ΠP of Π, the base of simple roots associated to this choice of positive roots, such that

{b, g−α }α∈ΠP generates p. In other words, parabolics containing a single Borel subgroup are

classified by subsets of the Dynkin diagram, with the empty set corresponding to the Borel,

and the whole graph corresponding to the group G.

Version: 1 Owner: bwebste Author(s): bwebste

278.16 pictures of Dynkin diagrams

**Here is a complete list of connected Dynkin diagrams. In general if the name of a diagram
**

has n as a subscript then there are n dots in the diagram. There are four infinite series that

correspond to classical complex (that is over C) simple Lie algebras. No pan intended.

• An , for n > 1 represents the simple complex Lie algebra sln+1 :

A1

A2

A3

An

• Bn , for n > 1 represents the simple complex Lie algebra so2n+1 :

• Cn , for n > 1 represents the simple complex Lie algebra sp2n :

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B1

B2

B3

Bn

C1

C2

C3

Cn

• Dn , for n > 3 represents the simple complex Lie algebra so2n :

D3

D4

D5

Dn

**And then there are the exceptional cases that come in finite families. The corresponding
**

Lie algebras are usually called by the name of the diagram.

**• There is the E series that has three members: E6 which represents a 78–dimensional Lie
**

algebra, E7 which represents a 133–dimensional Lie algebra, and E8 which represents

a 248–dimensional Lie algebra.

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E6

E7

E8

• There is the F4 diagram which represents a 52–dimensional complex simple Lie algebra:

F4

• And finally there is G2 that represents a 14–dimensional Lie algebra.

G2

**Notice the low dimensional coincidences:
**

A1 = B1 = C1

which reflects the exceptional isomorphisms

sl2 ∼

= so3 ∼

= sp2 .

Also

B2 ∼

= C2

reflecting the isomorphism

so5 ∼

= sp4 .

And,

A3 ∼

= D3

reflecting

sl4 ∼

= so6 .

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Remark 1. Often in the literature the listing of Dynkin diagrams is arranged so that there

are no “intersections” between different families. However by allowing intersections one gets

a graphical representation of the low degree isomorphisms. In the same vein there is a

graphical representation of the isomorphism

so4 ∼

= sl2 × sl2 .

Namely, if not for the requirement that the families consist of connected diagrams, one could

start the D family with

D2

which consists of two disjoint copies of A2 .

Version: 9 Owner: Dr Absentius Author(s): Dr Absentius

278.17 positive root

**If R ⊂ E is a root system, with E a euclidean vector space, then a subset R+ ⊂ R is called a
**

set of positive roots if there is a vector v ∈ E such that (α, v) > 0 if α ∈ R+ , and (α, v) < 0

if α ∈ R\R+ . roots which are not positive are called negative. Since −α is negative exactly

when α is positive, exactly half the roots must be positive.

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278.18 rank

**Let lg be a finite dimensional Lie algebra. One can show that all Cartan subalgebras h ⊂ lg
**

have the same dimension. The rank of lg is defined to be this dimension.

Version: 5 Owner: rmilson Author(s): rmilson

278.19 root lattice

**If R ⊂ E is a root system, and E a euclidean vector space, then the root lattice ΛR of R
**

is the subset of E generated by R as an abelian group. In fact, this group is free on the

simple roots, and is thus a full sublattice of E.

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278.20 root system

Root systems are sets of vectors in a Euclidean space which are used classify simple Lie algebras,

and to understand their representation theory, and also in the theory of reflection groups.

**Axiomatically, an (abstract) root system R is a set of vectors in a euclidean vector space E
**

with inner product (·, ·), such that:

**1. R spans the vector space E.
**

2. if α ∈ R, then reflection in the hyperplane orthogonal to α preserves R.

(α,β)

3. if α, β ∈ R, then 2 (α,α) is an integer.

**Axiom 3 is sometimes dropped when dealing with reflection groups, but it is necessary for
**

the root systems which arise in connection with Lie algebras.

Additionally, a root system is called reduced if for all α ∈ R, if kα ∈ R, then k = ±1.

**We call a root system indecomposable if there is no proper subset R0 ⊂ R such that every
**

vector in R0 is orthogonal to R.

**Root systems arise in the classification of semi-simple Lie algebras in the following manner:
**

If g is a semi-simple complex Lie algebra, then one can choose a maximal self-normalizing

subalgebra of g (alternatively, this is the commutant of an element with commutant of

minimal dimension), called a Cartan subalgebra, traditionally denote h. These act on g by

the adjoint action by diagonalizable linear maps. Since these maps all commute, they are all

simultaneously diagonalizable. The simultaneous eigenspaces of this action are called root

spaces, and the decomposition of g into h and the root spaces is called a root decompositon

of g. It turns out that all root spaces are all one dimensional. Now, for each eigenspace,

we have a map λ : h → C, given by Hv = λ(H)v for v an element of that eigenspace. The

set R ⊂ h∗ of these λ is called the root system of the algebra g. The Cartan subalgebra h

has a natural inner product (the Killing form), which in turn induces an inner product on

h∗ . With respect to this inner product, the root system R is an abstract root system, in the

sense defined up above.

**Conversely, given any abstract root system R, there is a unique semi-simple complex Lie
**

algebra g such that R is its root system. Thus to classify complex semi-simple Lie algebras,

we need only classify roots systems, a somewhat easier task. Really, we only need to classify

indecomposable root systems, since all other root systems are built out of these. The Lie

algebra corresponding to a root system is simple if and only if the associated root system is

indecomposable.

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By convention e1 , . . . , en are orthonormal vectors, and the subscript on the name of the root

system is the dimension of the space it is contained in, also called the rank of the system,

and the indices i and j will run from 1 to n. There are four infinite series of indecomposable

root systems :

P

• An = {ei − ej , δ + ei }i6=j , where δ = nk=1 ek . This system corresponds to sl2 C.

S

• Bn = {±ei ± ej }i<j {ei }. This system corresponds to so2n+1 C.

S

• Cn = {±ei ± ej }i<j {2ei }. This system corresponds to sp2n C.

• Dn = {±ei ± ej }i<j . This sytem corresponds to so2n C.

**and there are five exceptional root systems G2 , F4 , E6 , E7 , E8 , with five corresponding excep-
**

tional algebras, generally denoted by the same letter in lower-case Fraktur (g2 , etc.).

Version: 3 Owner: bwebste Author(s): bwebste

278.21 simple and semi-simple Lie algebras

**A Lie algebra is called simple if it has no proper ideals and is not abelian. A Lie algebra is
**

called semisimple if it has no proper solvable ideals and is not abelian.

**Let k = R or C. Examples of simple algebras are sln k, the Lie algebra of the special linear group
**

(traceless matrices), son k, the Lie algebra of the special orthogonal group (skew-symmetric matrices),

and sp2n k the Lie algebra of the symplectic group. Over R, there are other simple Lie al-

gebas, such as sun , the Lie algebra of the special unitary group (skew-Hermitian matrices).

Any semisimple Lie algebra is a direct product of simple Lie algebras.

Simple and semi-simple Lie algebras are one of the most widely studied classes of algebras for

a number of reasons. First of all, many of the most interesting Lie groups have semi-simple

Lie algebras. Secondly, their representation theory is very well understood. Finally, there is

a beautiful classification of simple Lie algebras.

Over C, there are 3 infinite series of simple Lie algebras: sln , son and sp2n , and 5 exceptional

simple Lie algebras g2 , f4 , e6 , e7 , and e8 . Over R the picture is more complicated, as several

different Lie algebras can have the same complexification (for example, sun and sln R both

have complexification sln C).

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278.22 simple root

**Let R ⊂ E be a root system, with E a euclidean vector space. If R+ is a set of positive roots,
**

then a root is called simple if it is positive, and not the sum of any two positive roots. The

simple roots form a basis of the vector space E, and any positive root is a positive integer

linear combination of simple roots.

**A set of roots which is simple with respect to some choice of a set of positive roots is called
**

a base. The Weyl group of the root system acts simply transitively on the set of bases.

Version: 1 Owner: bwebste Author(s): bwebste

278.23 weight (Lie algebras)

**Let g be a semi-simple Lie algebra. Choose a Cartan subalgebra h. Then a weight is simply
**

an element of the dual h∗ . Weights arise in the representation theory of semi-simple Lie

algebras in the following manner: The elements of h must act on V by diagonalizable (also

called semi-simple) linear transformations. Since h is abelian, these must be simultaneously

diagonalizable. Thus, V decomposes as the direct sum of simultaneous eigenspaces for h.

Let V be such an eigenspace. Then the map λ defined by λ(H)v = Hv is a linear functional

on h, and thus a weight, as defined above. The maximal eigenspace Vλ with weight λ is called

the weight space of λ. The dimension of Vλ is called the multiplicity of λ. A representation

of a semi-simple algebra is determine by the multiplicities of its weights.

Version: 3 Owner: bwebste Author(s): bwebste

278.24 weight lattice

**The weight lattice ΛW of a root system R ⊂ E is the dual lattice to ΛR , the root lattice of
**

R. That is,

ΛW = {e ∈ E|(e, r) ∈ Z}.

Weights which lie in the weight lattice are called integral. Since the simple roots are free

generators of the root lattice, one need only check that (e, π) ∈ Z for all simple roots π. If

R ⊂ h is the root system of a semi-simple Lie algebra g with Cartan subalgebra h, then ΛW

is exactly the set of weights appearing in finite dimensional representations of g.

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Chapter 279

**17B30 – Solvable, nilpotent
**

(super)algebras

279.1 Engel’s theorem

**Before proceeding, it will be useful to recall the definition of a nilpotent Lie algebra. Let g
**

be a Lie algebra. The lower central series of g is defined to be the filtration of ideals

D0 g ⊃ D1 g ⊃ D2 g ⊃ . . . ,

where

D0 g = g, Dk+1g = [g, Dk g], k ∈ N.

To say that g is nilpotent is to say that the lower central series has a trivial termination, i.e.

that there exists a k such that

Dk g = 0,

or equivalently, that k nested bracket operations always vanish.

**Theorem 1 (Engel). Let g ⊂ End V be a Lie algebra of endomorphisms of a finite-dimensional
**

vector space V . Suppose that all elements of g are nilpotent transformations. Then, g is a

nilpotent Lie algebra.

**Lemma 3. Let X : V → V be a nilpotent endomorphism of a vector space V . Then, the
**

adjoint action

ad(X) : End V → End V

is also a nilpotent endomorphism.

**Proof. Suppose that
**

Xk = 0

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for some k ∈ N. We will show that

ad(X)2k−1 = 0.

Note that

ad(X) = l(X) − r(X),

where

l(X), r(X) : End V → End V,

are the endomorphisms corresponding, respectively, to left and right multiplication by X.

These two endomorphisms commute, and hence we can use the binomial formula to write

2k−1

X

2k−1

ad(X) = (−1)i l(X)2k−1−i r(X)i .

i=0

Each of terms in the above sum vanishes because

l(X)k = r(X)k = 0.

QED

**Lemma 4. Let g be as in the theorem, and suppose, in addition, that g is a nilpotent Lie
**

algebra. Then the joint kernel, \

ker g = ker a,

a∈g

is non-trivial.

**Proof. We proceed by induction on the dimension of g. The claim is true for dimen-
**

sion 1, because then g is generated by a single nilpotent transformation, and all nilpotent

transformations are singular.

Suppose then that the claim is true for all Lie algebras of dimension less than n = dim g.

We note that D1 g fits the hypotheses of the lemma, and has dimension less than n, because

g is nilpotent. Hence, by the induction hypothesis

V0 = ker D1 g

**is non-trivial. Now, if we restrict all actions to V0 , we obtain a representation of g by abelian
**

transformations. This is because for all a, b ∈ g and v ∈ V0 we have

abv − bav = [a, b]v = 0.

**Now a finite number of mutually commuting linear endomorphisms admits a mutual eigenspace
**

decomposition. In particular, if all of the commuting endomorphisms are singular, their joint

kernel will be non-trivial. We apply this result to a basis of g/D1 g acting on V0 , and the

desired conclusion follows. QED

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Proof of the theorem. We proceed by induction on the dimension of g. The theorem is

true in dimension 1, because in that circumstance D1 g is trivial.

Next, suppose that the theorem holds for all Lie algebras of dimension less than n = dim g.

Let h ⊂ g be a properly contained subalgebra of minimum codimension. We claim that there

exists an a ∈ g but not in h such that [a, h] ⊂ h.

**By the induction hypothesis, h is nilpotent. To prove the claim consider the isotropy representation
**

of h on g/h. By Lemma 1, the action of each a ∈ h on g/h is a nilpotent endomorphism.

Hence, we can apply Lemma 2 to deduce that the joint kernel of all these actions is non-

trivial, i.e. there exists a a ∈ g but not in h such that

[b, a] ⇔ 0 (mod#1) ,

for all b ∈ h. Equivalently, [h, a] ⊂ h and the claim is proved.

**Evidently then, the span of a and h is a subalgebra of g. Since h has minimum codimension,
**

we infer that h and a span all of g, and that

D1 g ⊂ h. (279.1.1)

Next, we claim that all the Dk h are ideals of g. It is enough to show that

[a, Dk h] ⊂ Dk h.

**We argue by induction on k. Suppose the claim is true for some k. Let b ∈ h, c ∈ Dk h be
**

given. By the Jacobi identity

[a, [b, c]] = [[a, b], c] + [b, [a, c]].

**The first term on the right hand-side in Dk+1h because [a, b] ∈ h. The second term is in
**

Dk+1 h by the induction hypothesis. In this way the claim is established.

Now a is nilpotent, and hence by Lemma 1,

ad(a)n = 0 (279.1.2)

for some n ∈ N. We now claim that

Dn+1 g ⊂ D1 h.

**By (278.1.1) it suffices to show that
**

n times

z }| {

[g, [. . . [g, h] . . .]] ⊂ D1 h.

Putting

g1 = g/D1 h, h1 = h/D1 h,

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this is equivalent to

n times

z }| {

[g1 , [. . . [g1 , h1 ] . . .]] = 0.

However, h1 is abelian, and hence, the above follows directly from (278.1.2).

Adapting this argument in the obvious fashion we can show that

Dkn+1g ⊂ Dk h.

Since h is nilpotent, g must be nilpotent as well. QED

**Historical remark. In the traditional formulation of Engel’s theorem, the hypotheses are
**

the same, but the conclusion is that there exists a basis B of V , such that all elements of g

are represented by nilpotent matrices relative to B.

**Let us put this another way. The vector space of nilpotent matrices Nil, is a nilpotent Lie
**

algebra, and indeed all subalgebras of Nil are nilpotent Lie algebras. Engel’s theorem asserts

that the converse holds, i.e. if all elements of a Lie algebra g are nilpotent transformations,

then g is isomorphic to a subalgebra of Nil.

**The classical result follows straightforwardly from our version of the Theorem and from
**

Lemma 2. Indeed, let V1 be the joint kernel g. We then let U2 be the joint kernel of g acting

on V /V0 , and let V2 ⊂ V be the subspace obtained by pulling U2 x back to V . We do this a

finite number of times and obtain a flag of subspaces

0 = V0 ⊂ V1 ⊂ V2 ⊂ . . . ⊂ Vn = V,

such that

gVk+1 = Vk

for all k. The choose an adapted basis relative to this flag, and we’re done.

Version: 2 Owner: rmilson Author(s): rmilson

279.2 Lie’s theorem

**Let g be a finite dimensional complex solvable Lie algebra, and V a repesentation of g. Then
**

there exists an element of V which is a simultaneous eigenvector for all elements of g.

**Applying this result inductively, we find that there is a basis of V with respect to which all
**

elements of g are upper triangular.

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279.3 solvable Lie algebra

Let g be a Lie algebra. The lower central series of g is the filtration of subalgebras

D1 g ⊃ D2 g ⊃ D3 g ⊃ · · · ⊃ Dk g ⊃ · · ·

of g, inductively defined for every natural number k as follows:

D1 g := [g, g]

Dk g := [g, Dk−1 g]

The upper central series of g is the filtration

D1 g ⊃ D2 g ⊃ D3 g ⊃ · · · ⊃ Dk g ⊃ · · ·

defined inductively by

D1 g := [g, g]

Dk g := [Dk−1g, Dk−1 g]

**In fact both Dk g and Dk g are ideals of g, and Dk g ⊂ Dk g for all k. The Lie algebra g is
**

defined to be nilpotent if Dk g = 0 for some k ∈ N, and solvable if Dk g = 0 for some k ∈ N.

**A subalgebra h of g is said to be nilpotent or solvable if h is nilpotent or solvable when
**

considered as a Lie algebra in its own right. The terms may also be applied to ideals of g,

since every ideal of g is also a subalgebra.

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1183

Chapter 280

**17B35 – Universal enveloping
**

(super)algebras

280.1 Poincaré-Birkhoff-Witt theorem

**Let g be a Lie algebra over a field k, and let B be a k-basis of g equipped with a linear order
**

6. The Poincaré-Birkhoff-Witt-theorem (often abbreviated to PBW-theorem) states that

the monomials

x1 x2 · · · xn with x1 6 x2 6 . . . 6 xn elements of B

constitute a k-basis of the universal enveloping algebra U(g) of g. Such monomials are often

called ordered monomials or PBW-monomials.

It is easy to see that they span U(g): for all n ∈ N, let Mn denote the set

Mn = {(x1 , . . . , xn ) | x1 6 . . . 6 xn } ⊂ B n ,

S∞

and denote by π : n=0 B n → U(g) the multiplication map. Clearly it suffices to prove that

n

X

π(B n ) ⊆ π(Mi )

i=0

**for all n ∈ N; to this end, we proceed by induction. For n = 0 the statement is clear. Assume
**

that it holds for n − 1 > 0, and consider a list (x1 , . . . , xn ) ∈ B n . If it is an element of Mn ,

then we are done. Otherwise, there exists an index i such that xi > xi+1 . Now we have

π(x1 , . . . , xn ) = π(x1 , . . . , xi−1 , xi+1 , xi , xi+2 , . . . , xn )

+ x1 · · · xi−1 [xi , xi+1 ]xi+1 · · · xn .

As B is a basis of k, [xi , xi+1 ] is a P

linear combination of B. Using this to expand the second

n−1

term above, we find that it is in i=0 π(Mi ) by the induction hypothesis. The argument

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of π in the first term, on the other hand, is lexicographically smaller than (x1 , . . . , xn ), but

contains the same entries. Clearly this rewriting proces must end, and this concludes the

induction step.

The proof of linear independence of the PBW-monomials is slightly more difficult.

Version: 1 Owner: draisma Author(s): draisma

280.2 universal enveloping algebra

**A universal enveloping algebra of a Lie algebra g over a field k is an associative algebra U
**

(with unity) over k, together with a Lie algebra homomorphism ι : g → U (where the Lie

algebra structure on U is given by the commutator), such that if A is a another associative

algebra over k and φ : g → A is another Lie algebra homomorphism, then there exists a

unique homomorphism ψ : U → A of associative algebras such that the diagram

ι

g U

ψ

φ

A

commutes. Any g has a universal enveloping algebra: let T be the associative tensor algebra

generated by the vector space g, and let I be the two-sided ideal of T generated by elements

of the form

xy − yx − [x, y] for x, y ∈ g;

then U = T /I is a universal enveloping algebra of g. Moreover, the universal property above

ensures that all universal enveloping algebras of g are canonically isomorphic; this justifies

the standard notation U(g).

Some remarks:

**1. By the Poincaré-Birkhoff-Witt theorem, the map ι is injective; usually g is identified
**

with ι(g). From the construction above it is clear that this space generates U(g) as an

associative algebra with unity.

2. By definition, the (left) representation theory of U(g) is identical to that of g. In

particular, any irreducible g-module corresponds to a maximal left ideal of U(g).

**Example: let g be the Lie algebra generated by the elements p, q, and e with Lie bracket
**

determined by [p, q] = e and [p, e] = [q, e] = 0. Then U(g)/(e − 1) (where (e − 1) denotes the

∂

two-sided ideal generated by e − 1) is isomorphic to the skew polynomial algebra k[x, ∂x ],

the isomorphism being determined by

∂

p + (e − 1) 7→ and

∂x

q + (e − 1) 7→ x.

1185

Version: 1 Owner: draisma Author(s): draisma

1186

Chapter 281

**17B56 – Cohomology of Lie
**

(super)algebras

281.1 Lie algebra cohomology

Let g be a Lie algebra, and M a representation of g. Let

M g = {m ∈ M : Xm = 0∀X ∈ g}.

**This is clearly a covariant functor. Call its derived functor Ri (−g) = H i(g, −) the Lie algebra
**

cohomology of g with coefficients in M

**These cohomology groups have certain interpretations. For any Lie algebra, H 1 (g, k) ∼
**

=

2

g/[g, g], the abelianization of g, and H (g, M) is in natural bijection with Lie algebra

extensions (thinking of M as an abelian Lie algebra) 0 → M → f → g → 0 such that

the action of g on M induced by that of f coincides with that already specified.

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1187

Chapter 282

**17B67 – Kac-Moody (super)algebras
**

(structure and representation theory)

282.1 Kac-Moody algebra

**Let A be an n × n generalized Cartan matrix. If n − r is the rank of A, then let h be a n + r
**

dimensional complex vector space. Choose n linearly independent elements α1 , . . . , αn ∈ h∗

(called roots), and α̌1 , . . . , α̌n ∈ h (called coroots) such that hαi , α̌j i = aij , where h·, ·i is the

natural pairing of h∗ and h. This choice is unique up to automorphisms of h.

**Then the Kac-Moody algebra associated to g(A) is the Lie algebra generated by elements
**

X1 , . . . , Xn , Y1 , . . . , Yn and h, with the relations

[Xi , Yi ] = α̌i [Xi , Yj ] = 0

[Xi , h] = αi (h)Xi [Yi , h] = −αi (h)Yi

[X , [Xi , · · · , [Xi , Xj ] · · · ]] = 0 [Y , [Y , · · · , [Yi, Yj ] · · · ]] = 0

| i {z } | i i{z }

1−aij times 1−aij times

**If the matrix A is positive-definite, we obtain a finite dimensional semi-simple Lie algebra,
**

and A is the Cartan matrix associated to a Dynkin diagram. Otherwise, the algebra we

obtain is infinite dimensional and has an r-dimensional center.

Version: 2 Owner: bwebste Author(s): bwebste

282.2 generalized Cartan matrix

**A generalized Cartan matrix is a matrix A whose diagonal entries are all 2, and whose
**

off-diagonal entries are nonpositive integers, such that aij = 0 if and only if aji = 0. Such a

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matrix is called symmetrizable if there is a diagonal matrix B such that AB is symmetric.

Version: 2 Owner: bwebste Author(s): bwebste

1189

Chapter 283

17B99 – Miscellaneous

283.1 Jacobi identity interpretations

**The Jacobi identity in a Lie algebra g has various interpretations that are more transparent,
**

whence easier to remember, than the usual form

[x, [y, z]] + [y, [z, x]] + [z, [x, y]] = 0.

**One is the fact that the adjoint representation ad : g → End(g) really is a representation.
**

Yet another way to formulate the identity is

ad(x)[y, z] = [ad(x)y, z] + [y, ad(x)z],

i.e., ad(x) is a derivation on g for all x ∈ g.

Version: 2 Owner: draisma Author(s): draisma

283.2 Lie algebra

**A Lie algebra over a field k is a vector space g with a bilinear map [ , ] : g × g → g, called
**

the Lie bracket and denoted (x, y) 7→ [x, y]. It is required to satisfy:

1. [x, x] = 0 for all x ∈ g.

2. The Jacobi identity: [x, [y, z]] + [y, [z, x]] + [z, [x, y]] = 0 for all x, y, z ∈ g.

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283.2.1 Subalgebras & Ideals

**A vector subspace h of the Lie algebra g is a subalgebra if h is closed under the Lie bracket
**

operation, or, equivalently, if h itself is a Lie algebra under the same bracket operation as g.

An ideal of g is a subspace h for which [x, y] ∈ h whenever either x ∈ h or y ∈ h. Note that

every ideal is also a subalgebra.

Some general examples of subalgebras:

• The center of g, defined by Z(g) := {x ∈ g | [x, y] = 0for all y ∈ g}. It is an ideal of g.

**• The normalizer of a subalgebra h is the set N(h) := {x ∈ g | [x, h] ⊂ h}. The Jacobi
**

identity guarantees that N(h) is always a subalgebra of g.

**• The centralizer of a subset X ⊂ g is the set C(X) := {x ∈ g | [x, X] = 0}. Again, the
**

Jacobi identity implies that C(X) is a subalgebra of g.

283.2.2 Homomorphisms

**Given two Lie algebras g and g0 over the field k, a homomorphism from g to g0 is a
**

linear transformation φ : g → g0 such that φ([x, y]) = [φ(x), φ(y)] for all x, y ∈ g. An

injective homomorphism is called a monomorphism, and a surjective homomorphism is called

an epimorphism.

**The kernel of a homomorphism φ : g → g0 (considered as a linear transformation) is denoted
**

ker (φ). It is always an ideal in g.

283.2.3 Examples

• Any vector space can be made into a Lie algebra simply by setting [x, x] = 0 for all x.

The resulting Lie algebra is called an abelian Lie algebra.

**• If G is a Lie group, then the tangent space at the identity forms a Lie algebra over the
**

real numbers.

**• R3 with the cross product operation is a nonabelian three dimensional Lie algebra over
**

R.

283.2.4 Historical Note

**Lie algebras are so-named in honour of Sophus Lie, a Norwegian mathematician who pio-
**

neered the study of these mathematical objects. Lie’s discovery was tied to his investigation

1191

of continuous transformation groups and symmetries. One joint project with Felix Klein

called for the classification of all finite-dimensional groups acting on the plane. The task

seemed hopeless owing to the generally non-linear nature of such group actions. However,

Lie was able to solve the problem by remarking that a transformation group can be locally

reconstructed from its corresponding “infinitesimal generators”, that is to say vector fields

corresponding to various 1–parameter subgroups. In terms of this geometric correspondence,

the group composition operation manifests itself as the bracket of vector fields, and this is

very much a linear operation. Thus the task of classifying group actions in the plane became

the task of classifying all finite-dimensional Lie algebras of planar vector field; a project that

Lie brought to a successful conclusion.

**This “linearization trick” proved to be incredibly fruitful and led to great advances in
**

geometry and differential equations. Such advances are based, however, on various results

from the theory of Lie algebras. Lie was the first to make significant contributions to this

purely algebraic theory, but he was surely not the last.

Version: 10 Owner: djao Author(s): djao, rmilson, nerdy2

283.3 real form

**Let G be a complex Lie group. A real Lie group K called a real form of G if g ∼
**

= C ⊗R k,

where g and k are the Lie algebras of G and K, respectively.

Version: 2 Owner: bwebste Author(s): bwebste

1192

Chapter 284

**18-00 – General reference works
**

(handbooks, dictionaries,

bibliographies, etc.)

284.1 Grothendieck spectral sequence

**If F : C → D and G : D → E are two covariant left exact functors between abelian categories,
**

and if F takes injective objects of C to G-acyclic objects of D then there is a spectral sequence

for each object A of C:

E2pq = (Rp G ◦ Rq F )(A) → Rp+q (G ◦ F )(A)

**If X and Y are topological spaces and C = Ab(X) is the category of sheaves of abelian groups
**

on X and D = Ab(Y ) and E = Ab is the category of abelian groups, then for a continuous map

f : X → Y we have a functor f∗ : Ab(X) → Ab(Y ), the direct image functor. We also

have the global section functors ΓX : Ab(X) → Ab, and ΓY : Ab(Y ) → Ab. Then since

ΓY ◦ f∗ = ΓX and we can verify the hypothesis (injectives are flasque, direct images of

flasque sheaves are flasque, and flasque sheaves are acyclic for the global section functor),

the sequence in this case becomes:

H p (Y, Rq f∗ F) → H p+q (X, F)

for a sheaf F of abelian groups on X, exactly the Leray spectral sequence.

**I can recommend no better book than Weibel’s book on homological algebra. Sheaf theory
**

can be found in Hartshorne or in Godement’s book.

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284.2 category of sets

The category of sets has as its objects all sets and as its morphisms functions between sets.

(This works if a category’s objects are only required to be part of a class, as the class of all

sets exists.)

**Alternately one can specify a universe, containing all sets of interest in the situation, and
**

take the category to contain only sets in that universe and functions between those sets.

Version: 1 Owner: nerdy2 Author(s): nerdy2

284.3 functor

**Given two categories C and D, a covariant functor T : C → D consists of an assignment
**

for each object X of C an object T (X) of D (i.e. a “function” T : Ob(C) → Ob(D))

together with an assignment for every morphism f ∈ HomC(A, B), to a morphism T (f ) ∈

HomD(T (A), T (B)), such that:

**• T (1A ) = 1T (A) where 1X denotes the identity morphism on the object X (in the re-
**

spective category).

• T (g ◦ f ) = T (g) ◦ T (f ), whenever the composition g ◦ f is defined.

**A contravariant functor T : C → D is just a covariant functor T : Cop → D from the
**

opposite category. In other words, the assignment reverses the direction of maps. If f ∈

HomC(A, B), then T (f ) ∈ HomD(T (B), T (A)) and T (g ◦ f ) = T (f ) ◦ T (g) whenever the

composition is defined (the domain of g is the same as the codomain of f ).

**Given a category C and an object X we always have the functor T : C → Sets to the
**

category of sets defined on objects by T (A) = Hom(X, A). If f : A → B is a morphism of C,

then we define T (f ) : Hom(X, A) → Hom(X, B) by g 7→ f ◦ g. This is a covariant functor,

denoted by Hom(X, −).

Similarly, one can define a contravariant functor Hom(−, X) : C → Sets.

Version: 3 Owner: nerdy2 Author(s): nerdy2

284.4 monic

**A morphism f : A → B in a category is called monic if for any object C and any morphisms
**

g1 , g2 : C → A, if f ◦ g1 = f ◦ g2 then g1 = g2 .

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A monic in the category of sets is simply a one-to-one function.

Version: 1 Owner: nerdy2 Author(s): nerdy2

284.5 natural equivalence

**A natural transformation between functors τ : F → G is called a natural equivalence (or a
**

natural isomorphism) if there is a natural transformation σ : G → F such that τ ◦ σ = idG

and σ ◦ τ = idF where idF is the identity natural transformation on F (which for each object

A gives the identity map F (A) → F (A)), and composition is defined in the obvious way

(for each object compose the morphisms and it’s easy to see that this results in a natural

transformation).

Version: 2 Owner: mathcam Author(s): mathcam, nerdy2

284.6 representable functor

**A contravariant functor T : C → Sets between a category and the category of sets is rep-
**

resentable if there is an object X of C such that T is isomorphic to the functor X • =

Hom(−, X).

Similarly, a covariant functor is T called representable if it is isomorphic to X• = Hom(X, −).

We say that the object X represents T . X is unique up to canonical isomorphism.

**A vast number of important objects in mathematics are defined as representing functors.
**

For example, if F : C → D is any functor, then the adjoint G : D → C (if it exists)

can be defined as follows. For Y in D, G(Y ) is the object of C representing the functor

X 7→ Hom(F (X), Y ) if G is right adjoint to F or X 7→ Hom(Y, F (X)) if G is left adjoint.

Thus, for example, if R is a ring, then N⊗M represents the functor L 7→ HomR (N, HomR (M, L)).

Version: 3 Owner: bwebste Author(s): bwebste, nerdy2

284.7 supplemental axioms for an Abelian category

**These are axioms introduced by Alexandre Grothendieck for an abelian category. The first
**

two are satisfied by definition in an Abelian category, and others may or may not be.

(Ab1) Every morphism has a kernel and a cokernel.

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(Ab2) Every monic is the kernel of its cokernel.

**(Ab3) coproducts exist. (Coproducts are also called direct sums.) If this axiom is satisfied
**

the category is often just called cocomplete.

(Ab3*) Products exist. If this axiom is satisfied the category is often just called complete.

(Ab4) Coproducts exist and the coproduct of monics is a monic.

(Ab4*) Products exist and the product of epics is an epic.

(Ab5) Coproducts exist and filtered colimits of exact sequences are exact.

(Ab5*) Products exist and filtered inverse limits of exact sequences are exact.

**Grothendieck introduced these in his homological algebra paper in the Tokohu Journal of
**

Math. They can also be found in Weibel’s excellent homological algebra book.

Version: 5 Owner: nerdy2 Author(s): nerdy2

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Chapter 285

18A05 – Definitions, generalizations

285.1 autofunctor

**Let F : C → C be an endofunctor on a category C. If F is a bijection on objects, Ob(C), and
**

morphisms, Mor(C), then it is an autofunctor.

**In short, an autofunctor is a full and faithful endofunctor F : C → C such that the mapping
**

Ob(C) → Ob(C) which is induced by F is a bijection.

An autofunctor F : C → C is naturally isomorphic to the identity functor idC.

Version: 10 Owner: mathcam Author(s): mathcam, mhale, yark, gorun manolescu

285.2 automorphism

**Roughly, an automorphism is a map from a mathematical object onto itself such that: 1.
**

There exists an ”inverse” map such that the composition of the two is the identity map of

the object, and 2. any relevent structure related to the object in question is preserved.

**In category theory, an automorphism of an object A in a category C is a morphishm ψ ∈
**

Mor(A, A) such that there exists another morphism φ ∈ Mor(A, A) and ψ ◦ φ = φ ◦ ψ = idA .

**For example in the category of groups an automorphism is just a bijective (inverse exists and
**

composition gives the identity) group homomorphism (group structure is preserved). Con-

cretely, the map: x 7→ −x is an automorphism of the additive group of real numbers. In the

category of topological spaces an automorphism would be a bijective, continuous map such

that it’s inverse map is also continuous (not guaranteed as in the group case). Concretely,

the map ψ : S 1 → S 1 where ψ(α) = α + θ for some fixed angle θ is an automorphism of the

topological space that is the circle.

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Version: 4 Owner: benjaminfjones Author(s): benjaminfjones

285.3 category

A category C consists of the following data:

1. a collection ob(C) of objects (of C)

**2. for each ordered pair (A, B) of objects of C, a collection (we will assume it is a set)
**

Hom(A, B) of morphisms from the domain A to the codomain B

3. a function ◦ : Hom(A, B) × Hom(B, C) → Hom(A, C) called composition.

**We normally denote ◦(f, g) by g ◦ f for morphisms f, g. The above data must satisfy the
**

following axioms: for objects A, B, C, D,

T

A1: Hom(A, B) Hom(C, D) = ∅ whenever A 6= B or C 6= D

**A2: (associativity) if f ∈ Hom(A, B), g ∈ Hom(B, C) and h ∈ Hom(C, D), h ◦ (g ◦ f ) =
**

(h ◦ g) ◦ f

**A3: (Existence of an identity morphism) for each object A there exists an identity morphism
**

idA ∈ Hom(A, A) such that for every f ∈ Hom(A, B), f ◦ idA = f and idA ◦ g = g for every

g ∈ Hom(B, A).

Some examples of categories:

**• 0 is the empty category with no objects or morphisms, 1 is the category with one
**

object and one (identity) morphism.

**• If we assume we have a universe U which contains all sets encountered in “everyday”
**

mathematics, Set is the category of all such small sets with morphisms being set

functions

• Top is the category of all small topological spaces with morphisms continuous functions

• Grp is the category of all small groups whose morphisms are group homomorphisms

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1198

285.4 category example (arrow category)

**Let C be a category, and let D be the category whose objects are the arrows of C. A
**

morphism between two morphisms f : A → B and g : A0 → B 0 is defined to be a couple

of morphisms (h, k), where h ∈ Hom(A, A0 ) and k ∈ Hom(B, B 0 ) such that the following

diagram

h

A A0

f g

B k

B0

commutes. The resulting category D is called the arrow category of C.

Version: 6 Owner: n3o Author(s): n3o

285.5 commutative diagram

**Definition 15. Let C be a category. A diagram in C is a directed graph Γ with vertex
**

set V and edge set E, (“loops” and “parallel edges” are allowed) together with two maps

o : V → Obj(C), m : E → Morph(C) such that if e ∈ E has source s(e) ∈ V and target

t(e) ∈ V then m(e) ∈ HomC (s(e), t(e)).

**Usually diagrams are denoted by drawing the corresponding graph and labeling its vertices
**

(respectively edges) with their images under o (respectively m), for example if f : A → B is

a morphism

f

A B

is a diagram. Often (as in the previous example) the vertices themselves are not drawn since

their position can b deduced by the position of their labels.

Definition 16. Let D = (Γ, o, m) be a diagram in the category C and γ = (e1 , . . . , en ) be a

path in Γ. Then the composition along γ is the following morphism of C

◦(γ) := m(en ) ◦ · · · ◦ m(e1 ) .

We say that D is commutative or that it commutes if for any two objects in the image of

o, say A = o(v1 ) and B = o(v2 ), and any two paths γ1 and γ2 that connect v1 to v2 we have

◦(γ1 ) = ◦(γ2 ) .

**For example the commutativity of the triangle
**

f

A B

h g

C

1199

translates to h = g ◦ f , while the commutativity of the square

f

A B

k g

C h

D

translates to g ◦ f = h ◦ k.

Version: 3 Owner: Dr Absentius Author(s): Dr Absentius

285.6 double dual embedding

**Let V be a vector space over a field K. Recall that V ∗ , the dual space, is defined to be the
**

vector space of all linear forms on V . There is a natural embedding of V into V ∗∗ , the dual

of its dual space. In the language of categories, this embedding is a natural transformation

between the identity functor and the double dual functor, both endofunctors operating on

VK , the category of vector spaces over K.

**Turning to the details, let
**

I, D : VK → VK

denote the identity and the dual functors, respectively. Recall that for a linear mapping

L : U → V (a morphism in VK ), the dual homomorphism D[L] : V ∗ → U ∗ is defined by

D[L](α) : u 7→ α(Lu), u ∈ U, α ∈ V ∗ .

The double dual embedding is a natural transformation

δ : I → D2 ,

that associates to every V ∈ VK a linear homomorphism δV ∈ Hom(V, V ∗∗ ) described by

δV (v) : α 7→ α(v), v ∈ V, α ∈ V ∗

**To show that this transformation is natural, let L : U → V be a linear mapping. We must
**

show that the following diagram commutes:

δU

U U ∗∗

L D 2 [L]

δV

V V ∗∗

Let u ∈ U and α ∈ V ∗ be given. Following the arrows down and right we have that

(δV ◦ L)(u) : α 7→ α(Lu).

1200

Following the arrows right, then down we have that

(D[D[L]] ◦ δU )(u) : α 7→ (δU u)(D[L]α)

= (D[L]α)(u)

= α(Lu),

as desired.

**Let us also note that for every non-zero v ∈ V , there exists an α ∈ V ∗ such that α(v) 6= 0.
**

Hence δV (v) 6= 0, and hence δV is an embedding, i.e. it is one-to-one. If V is finite dimensional,

then V ∗ has the same dimension as V . Consequently, for finite-dimensional V , the natural

embedding δV is, in fact, an isomorphism.

Version: 1 Owner: rmilson Author(s): rmilson

285.7 dual category

**Let C be a category. The dual category C∗ of C is the category which has the same objects
**

as C, but in which all morphisms are ”reversed”. That is to say if A, B are objects of C and

we have a morphism f : A → B, then f ∗ : B → A is a morphism in C∗ . The dual category

is sometimes called the opposite category and is denoted Cop .

Version: 3 Owner: RevBobo Author(s): RevBobo

285.8 duality principle

**Let Σ be any statement of the elementary theory of an abstract category. We form the dual
**

of Σ as follows:

1. Replace each occurrence of ”domain” in Σ with ”codomain” and conversely.

2. Replace each occurrence of g ◦ f = h with f ◦ g = h

**Informally, these conditions state that the dual of a statement is formed by reversing arrows
**

and compositions. For example, consider the following statements about a category C:

• f :A→B

**• f is monic, i.e. for all morphisms g, h for which composition makes sense, f ◦ g = f ◦ h
**

implies g = h.

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The respective dual statements are

• f :B→A

• f is epi, i.e. for all morphisms g, h for which composition makes sense, g ◦ f = h ◦ f

implies g = h.

**The duality principle asserts that if a statement is a theorem, then the dual statment is
**

also a theorem. We take ”theorem” here to mean provable from the axioms of the elementary

theory of an abstract category. In practice, for a valid statement about a particular category

C, the dual statement is valid in the dual category C∗ (Cop ).

Version: 3 Owner: RevBobo Author(s): RevBobo

285.9 endofunctor

Given a category C, an endofunctor is a functor T : C → C.

Version: 2 Owner: rmilson Author(s): NeuRet, Logan

**285.10 examples of initial objects, terminal objects and
**

zero objects

Examples of initial objects, terminal objects and zero objects of categories include:

**• The empty set is the unique initial object in the category of sets; every one-element
**

set is a terminal object in this category; there are no zero objects. Similarly, the empty

space is the unique initial object in the category of topological spaces; every one-point

space is a terminal object in this category.

**• In the category of non-empty sets, there are no initial objects. The singletons are not
**

initial: while every non-empty set admits a function from a singleton, this function is

in general not unique.

**• In the category of pointed sets (whose objects are non-empty sets together with a distin-
**

guished point; a morphism from (A, a) to (B, b) is a function f : A → B with f (a) = b)

every singleton serves as a zero object. Similarly, in the category of pointed topological spaces,

every singleton is a zero object.

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• In the category of groups, any trivial group (consisting only of its identity element) is

a zero object. The same is true for the category of abelian groups as well as for the

category of modules over a fixed ring. This is the origin of the term ”zero object”.

**• In the category of rings with identity, the ring of integers (and any ring isomorphic to
**

it) serves as an initial object. The trivial ring consisting only of a single element 0 = 1

is a terminal object.

**• In the category of schemes, the prime spectrum of the integers spec(Z) is a terminal
**

object. The emtpy scheme (which is the prime spectrum of the trivial ring) is an initial

object.

• In the category of fields, there are no initial or terminal objects.

**• Any partially ordered set (P, ≤) can be interpreted as a category: the objects are the
**

elements of P , and there is a single morphism from x to y if and only if x ≤ y. This

category has an initial object if and only if P has a smallest element; it has a terminal

object if and only if P has a largest element. This explains the terminology.

**• In the category of graphs, the null graph is an initial object. There are no terminal
**

objects, unless we allow our graphs to have loops (edges starting and ending at the

same vertex), in which case the one-point-one-loop graph is terminal.

**• Similarly, the category of all small categories with functors as morphisms has the empty
**

category as initial object and the one-object-one-morphism category as terminal object.

**• Any topological space X can be viewed as a category X̂ by taking the open sets as
**

objects, and a single morphism between two open sets U and V if and only if U ⊂ V .

The empty set is the initial object of this category, and X is the terminal object.

**• If X is a topological space and C is some small category, we can form the category of
**

all contravariant functors from X̂ to C, using natural transformations as morphisms.

This category is called the category of presheaves on X with values in C. If C

has an initial object c, then the constant functor which sends every open set to c is an

initial object in the category of presheaves. Similarly, if C has a terminal object, then

the corresponding constant functor serves as a terminal presheave.

**• If we fix a homomorphism f : A → B of abelian groups, we can consider the category
**

C consisting of all pairs (X, φ) where X is an abelian group and φ : X → A is a

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group homomorphism with f φ = 0. A morphism from the pair (X, φ) to the pair

(Y, ψ) is defined to be a group homomorphism r : X → Y with the property ψr = φ:

X φ

f

r A B

Y ψ

**The kernel of f is a terminal object in this category; this expresses the universal prop-
**

erty of kernels. With an analogous construction, cokernels can be retrieved as initial

objects of a suitable category.

**• The previous example can be generalized to arbitrary limits of functors: if F : I → C is
**

a functor, we define a new category F̂ as follows: its objects are pairs (X, (φi )) where

X is an object of C and for every object i of I, φi : X → F (i) is a morphism in C such

that for every morphism ρ : i → j in I, we have F (ρ)φi = φj . A morphism between

pairs (X, (φi )) and (Y, (ψi )) is defined to be a morphism r : X → Y such that ψi r = φi

for all objects i of I. The universal property of the limit can then be expressed as

saying: any terminal object of F̂ is a limit of F and vice versa (note that F̂ need not

contain a terminal object, just like F need not have a limit).

Version: 11 Owner: AxelBoldt Author(s): AxelBoldt

285.11 forgetful functor

**Let C and D be categories such that each object c of C can be regarded an object of D
**

by suitably ignoring structures c may have as a C-object but not a D-object. A functor

U : C → D which operates on objects of C by “forgetting” any imposed mathematical

structure is called a forgetful functor. The following are examples of forgetful functors:

**1. U : Grp → Set takes groups into their underlying sets and group homomorphisms to
**

set maps.

**2. U : Top → Set takes topological spaces into their underlying sets and continuous maps
**

to set maps.

3. U : Ab → Grp takes abelian groups to groups and acts as identity on arrows.

Forgetful functors are often instrumental in studying adjoint functors.

Version: 1 Owner: RevBobo Author(s): RevBobo

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285.12 isomorphism

**A morphism f : A −→ B in a category is an isomorphism if there exists a morphism
**

f −1 : B −→ A which is its inverse. The objects A and B are isomorphic if there is an

isomorphism between them.

Examples:

**• In the category of sets and functions, a function f : A −→ B is an isomorphism if and
**

only if it is bijective.

**• In the category of groups and group homomorphisms (or rings and ring homomorphisms),
**

a homomorphism φ : G −→ H is an isomorphism if it has an inverse map φ−1 : H −→ G

which is also a homomorphism.

**• In the category of vector spaces and linear transformations, a linear transformation is
**

an isomorphism if and only if it is an invertible linear transformation.

**• In the category of topological spaces and continuous maps, a continuous map is an
**

isomorphism if and only if it is a homeomorphism.

Version: 2 Owner: djao Author(s): djao

285.13 natural transformation

**Let A, B be categories and T, S : A → B functors. A natural transformation τ : S → T
**

is a family of morphisms τ = {τA : T (A) → S(A)} such that for each object A of A,

0

τA : S(A) → T (A) is an object of B and for each morphism f : A → A in A the following

diagram commutes:

τA

S(A) T (A)

Sf Tf

0

τA0 0

S(A ) T (A )

Version: 6 Owner: RevBobo Author(s): RevBobo

285.14 types of homomorphisms

**Often in a category of algebraic structures, those structures are generated by certain ele-
**

ments, and subject to certain relations. One often refers to functions between structures

1205

which are said to preserve those relations. These functions are typically called homomor-

phisms.

**An example is the category of groups. Suppose that f : A → B is a function between two
**

groups. We say that f is a group homomorphism if:

(a) the binary operator is preserved: f (a1 · a2 ) = f (a1 ) · f (a2 ) for all a1 , a2 ∈ A;

(b) the identity element is preserved: f (eA ) = eB ;

(c) inverses of elements are preserved: f (a−1 ) = [f (a)]−1 for all a ∈ A.

**One can define similar natural concepts of homomorphisms for other algebraic structures,
**

giving us ring homomorphisms, module homomorphisms, and a host of others.

We give special names to homomorphisms when their functions have interesting properties.

**If a homomorphism is an injective function (i.e. one-to-one), then we say that it is a
**

monomorphism. These are typically monic in their category.

**If a homomorphism is an surjective function (i.e. onto), then we say that it is an epimorphism.
**

These are typically epic in their category.

**If a homomorphism is an bijective function (i.e. both one-to-one and onto), then we say that
**

it is an isomorphism.

**If the domain of a homomorphism is the same as its codomain (e.g. a homomorphism
**

f : A → A), then we say that it is an endomorphism. We often denote the collection of

endomorphisms on A as End(A).

**If a homomorphism is both an endomorphism and an isomorphism, then we say that it is an
**

automorphism. We often denote the collection of automorphisms on A as Aut(A).

Version: 4 Owner: antizeus Author(s): antizeus

285.15 zero object

**An initial object in a category C is an object A in C such that, for every object X in C, there
**

is exactly one morphism A −→ X.

**A terminal object in a category C is an object B in C such that, for every object X in C,
**

there is exactly one morphism X −→ B.

**A zero object in a category C is an object 0 that is both an initial object and a terminal
**

object.

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All initial objects (respectively, terminal objects, and zero objects), if they exist, are isomorphic

in C.

Version: 2 Owner: djao Author(s): djao

1207

Chapter 286

**18A22 – Special properties of functors
**

(faithful, full, etc.)

286.1 exact functor

**A covariant functor F is said to be left exact if whenever
**

α β

0 → A −→ B −→ C

**is an exact sequence, then
**

Fα Fβ

0 → F A −→ F B −→ F C

is also an exact sequence.

**A covariant functor F is said to be right exact if whenever
**

α β

A −→ B −→ C → 0

**is an exact sequence, then
**

Fα Fβ

F A −→ F B −→ F C → 0

is also an exact sequence.

**A contravariant functor F is said to be left exact if whenever
**

α β

A −→ B −→ C → 0

**is an exact sequence, then
**

Fβ Fα

0 → F C −→ F B −→ F A

is also an exact sequence.

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A contravariant functor F is said to be right exact if whenever

α β

0 → A −→ B −→ C

**is an exact sequence, then
**

Fβ Fα

F C −→ F B −→ F A → 0

is also an exact sequence.

**A (covariant or contravariant) functor is said to be exact if it is both left exact and right
**

exact.

Version: 3 Owner: antizeus Author(s): antizeus

1209

Chapter 287

**18A25 – Functor categories, comma
**

categories

287.1 Yoneda embedding

**If C is a category, write Ĉ for the category of contravariant functors from C to Sets, the
**

category of sets. The morphisms in Ĉ are natural transformations of functors.

**(To avoid set theoretical concerns, one can take a universe U and take all categories to be
**

U-small.)

**For any object X of C, there is the functor hX = Hom(−, X). Then X 7→ hX is a covariant
**

functor C → Ĉ, which embeds C faithfully as a full subcategory of Ĉ.

Version: 4 Owner: nerdy2 Author(s): nerdy2

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Chapter 288

**18A30 – Limits and colimits
**

(products, sums, directed limits,

pushouts, fiber products, equalizers,

kernels, ends and coends, etc.)

288.1 categorical direct product

**Let {Ci}i∈IQbe a set of objects in a category Q
**

C. A direct product of the collection {Ci }i∈I is

an object i∈I Ci of C, with morphisms πi : j∈I Cj −→ Ci for each i ∈ I, such that:

**For every object A in C, and any collection
**

Q of morphisms fi : A −→ Ci for every i ∈ I, there

exists a unique morphism f : A −→ i∈I Ci making the following diagram commute for all

i ∈ I.

fi

A Ci

f πi

Q

j∈I Cj

Version: 4 Owner: djao Author(s): djao

288.2 categorical direct sum

Let {C`

i }i∈I be a set of objects in a category C. `

A direct sum of the collection {Ci }i∈I is an

object i∈I Ci of C, with morphisms ιi : Ci −→ j∈I Cj for each i ∈ I, such that:

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` collection of morphisms fi : Ci −→ A for every i ∈ I, there

For every object A in C, and any

exists a unique morphism f : i∈I Ci −→ A making the following diagram commute for all

i ∈ I.

fi

Ci A

ιi f

`

j∈I Cj

Version: 4 Owner: djao Author(s): djao

288.3 kernel

**Let f : X → Y be a function and let Y be have some sort of zero, neutral or null element
**

that we’ll denote as e. (Examples are groups, vector spaces, modules, etc)

**The kernel of f is the set:
**

ker f = {x ∈ X : f (x) = e}

that is, the set of elements in X such that their image is e. This set can also denoted as

f −1 (e) (that doesn’t mean f has an inverse function, it’s just notation) and that is read

as ”the kernel is the preimage of the neutral element”. Let’s see an examples. If X = Z

and Y = Z6 , the function f that sends each integer n to its residue class modulo 6. So

f (4) = 4, f (20) = 2, f (−5) = 1. The kernel of f consist precisely of the multiples of 6 (since

they have residue 0, we have f (6k) = 0).

**This is also an example of kernel of a group homomorphism, and since the sets are also rings,
**

the function f is also a homomorphism between rings and the kernel is also the kernel of a

ring homomorphism.

**Usually we are interested on sets with certain algebraic structure. In particular, the following
**

theorem holds for maps between pairs of vector spaces, groups, rings and fields (and some

other algebraic structures):

A map f : X → Y is injective if and only if ker f = {0} (the zero of Y ).

Version: 4 Owner: drini Author(s): drini

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Chapter 289

**18A40 – Adjoint functors (universal
**

constructions, reflective

subcategories, Kan extensions, etc.)

289.1 adjoint functor

**Let C, D be categories and T : C → D, S : D → C be covariant functors. T is said to be a
**

left adjoint functor to S (equivalently, S is a right adjoint functor to T ) if there exists

ν = νC,D such that

ν : Hom(T (C), D) ∼

= Hom(C, S(D))

D C

is a natural bijection of hom-sets for all objects C of C and D of D.

An adjoint to any functor are unique up to natural transformation.

Examples:

**1. Let U : Top → Set be the forgetful functor (i.e. U takes topological spaces to their
**

underlying sets, and continuous maps to set functions). Then U is right adjoint to the

functor F : Set → Top which gives each set the discrete topology.

**2. If U : Grp → Set is again the forgetful functor, this time on the category of groups,
**

the functor F : Set → Grp which takes a set A to the free group generated by A is

left adjoint to U.

**3. If UN : R − mod → R − mod is the functor M 7→ N ⊗ M for an R-module N, then UN
**

is the left adjoint to the functor FN : R−mod → R−mod given by L 7→ HomR (N, L).

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Version: 8 Owner: bwebste Author(s): bwebste, RevBobo

289.2 equivalence of categories

**Let C and D be two categories with functors F : C → D and G : D → C. The functors F
**

and G are an

**Definition 17. equivalence of categories if there are natural isomorphisms F G ∼
**

= idD and

∼

GF = idC .

**Note, F is left adjoint to G, and G is right adjoint to F as
**

G

Hom(F (c), d) −→ Hom(GF (c), G(d)) ←→ Hom(c, G(d)).

D C C

**And, F is right adjoint to G, and G is left adjoint to F as
**

F

Hom(G(d), c) −→ Hom(F G(d), F (c)) ←→ Hom(d, F (c)).

C D D

**In practical terms, two categories are equivalent if there is a fully faithful functor F : C → D,
**

such that every object d ∈ D is isomorphic to an object F (c), for some c ∈ C.

Version: 2 Owner: mhale Author(s): mhale

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Chapter 290

**18B40 – Groupoids, semigroupoids,
**

semigroups, groups (viewed as

categories)

290.1 groupoid (category theoretic)

**A groupoid, also known as a virtual group, is a small category where every morphism is
**

invertible.

There is also a group-theoretic concept with the same name.

Version: 6 Owner: akrowne Author(s): akrowne

1215

Chapter 291

**18E10 – Exact categories, abelian
**

categories

291.1 abelian category

**An abelian category is a category A satisfying the following axioms. Because the later axioms
**

rely on terms whose definitions involve the earlier axioms, we will intersperse the statements

of the axioms with such auxiliary definitions as needed.

Axiom 1. For any two objects A, B in A, the set of morphisms Hom(A, B) is an abelian group.

**The identity element in the group Hom(·, ·) will be denoted by 0, and the group operation
**

by +.

**Axiom 2. Composition of morphisms distributes over addition in Hom(·, ·). That is, given
**

any diagram of morphisms

g1

f h

A B C D

g2

we have (g1 + g2 )f = g1 f + g2 f and h(g1 + g2 ) = hg1 + hg2 .

Axiom 3. A has a zero object.

Axiom 4. For any two objects A, B in A, the categorical direct product A × B exists in A.

Given a morphism f : A −→ B in A, a kernel of f is a morphism i : X −→ A such that:

• f i = 0.

• For any other morphism j : X 0 −→ A such that f j = 0, there exists a unique morphism

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j 0 : X 0 −→ X such that the diagram

X0

j0

j

i f

X A B

commutes.

Likewise, a cokernel of f is a morphism p : B −→ Y such that:

• pf = 0.

**• For any other morphism j : B −→ Y 0 such that jf = 0, there exists a unique morphism
**

j 0 : Y −→ Y 0 such that the diagram

f p

A B Y

j

j0

0

Y

commutes.

Axiom 5. Every morphism in A has a kernel and a cokernel.

The kernel and cokernel of a morphism f in A will be denoted ker (f ) and cok(f ), respectively.

**A morphism f : A −→ B in A is called a monomorphism if, for every morphism g : X −→ A
**

such that f g = 0, we have g = 0. Similarly, the morphism f is called an epimorphism if, for

every morphism h : B −→ Y such that hf = 0, we have h = 0.

Axiom 6. ker (cok(f )) = f for every monomorphism f in A.

Axiom 7. cok(ker (f )) = f for every epimorphism f in A.

Version: 6 Owner: djao Author(s): djao

291.2 exact sequence

Let A be an abelian category. We begin with a preliminary definition.

**Definition 1. For any morphism f : A −→ B in A, let m : X −→ B be the morphism
**

equal to ker (cok(f )). Then the object X is called the image of f , and denoted Im(f ). The

morphism m is called the image morphism of f , and denoted i(f ).

1217

Note that Im(f ) is not the same as i(f ): the former is an object of A, while the latter is a

morphism of A. We note that f factors through i(f ):

e i(f )

A Im(f ) B

f

**The proof is as follows: by definition of cokernel, cok(f )f = 0; therefore by definition of
**

kernel, the morphism f factors through ker (cok(f )) = i(f ), and this factor is the morphism

e above. Furthermore m is a monomorphism and e is an epimorphism, although we do not

prove these facts.

Definition 2. A sequence

f g

···A B C ···

of morphisms in A is exact at B if ker (g) = i(f ).

Version: 3 Owner: djao Author(s): djao

291.3 derived category

**Let A be an abelian category, and let K(A) be the category of chain complexes in A, with
**

morphisms chain homotopy classes of maps. Call a morphism of chain complexes a quasi-

isomorphism if it induces an isomorphism on homology groups of the complexes. For exam-

ple, any chain homotopy is an quasi-isomorphism, but not conversely. Now let the derived

category D(A) be the category obtained from K(A) by adding a formal inverse to every

quasi-isomorphism (technically this called a localization of the category).

**Derived categories seem somewhat obscure, but in fact, many mathematicians believe they
**

are the appropriate place to do homological algebra. One of their great advantages is that the

important functors of homological algebra which are left or right exact (Hom,N⊗k , where

N is a fixed k-module, the global section functor Γ, etc.) become exact on the level of

derived functors (with an appropriately modified definition of exact).

See Methods of Homological Algebra, by Gelfand and Manin for more details.

Version: 2 Owner: bwebste Author(s): bwebste

291.4 enough injectives

**An abelian category is said to have enough injectives if for every object X, there is a
**

monomorphism 0 → X → I where I is an injective object.

Version: 2 Owner: bwebste Author(s): bwebste

1218

Chapter 292

18F20 – Presheaves and sheaves

292.1 locally ringed space

292.1.1 Definitions

**A locally ringed space is a topological space X together with a sheaf of rings OX with the
**

property that, for every point p ∈ X, the stalk (OX )p is a local ring 1 .

**A morphism of locally ringed spaces from (X, OX ) to (Y, OY ) is a continuous map f : X −→
**

Y together with a morphism of sheaves φ : OY −→ OX with respect to f such that, for every

point p ∈ X, the induced ring homomorphism on stalks φp : (OY )f (p) −→ (OX )p is a local

homomorphism. That is,

φp (y) ∈ mp for every y ∈ mf (p) ,

where mp (respectively, mf (p) ) is the maximal ideal of the ring (OX )p (respectively, (OY )f (p) ).

292.1.2 Applications

**Locally ringed spaces are encountered in many natural contexts. Basically, every sheaf on
**

the topological space X consisting of continuous functions with values in some field is a

locally ringed space. Indeed, any such function which is not zero at a point p ∈ X is

nonzero and thus invertible in some neighborhood of p, which implies that the only maximal

ideal of the stalk at p is the set of germs of functions which vanish at p. The utility of

this definition lies in the fact that one can then form constructions in familiar instances of

locally ringed spaces which readily generalize in ways that would not necessarily be obvious

without this framework. For example, given a manifold X and its locally ringed space DX

of real–valued differentiable functions, one can show that the space of all tangent vectors to

1

All rings mentioned in this article are required to be commutative.

1219

X at p is naturally isomorphic to the real vector space (mp /m2p )∗ , where the ∗ indicates the

dual vector space. We then see that, in general, for any locally ringed space X, the space

of tangent vectors at p should be defined as the k–vector space (mp /m2p )∗ , where k is the

residue field (OX )p /mp and ∗ denotes dual with respect to k as before. It turns out that

this definition is the correct definition even in esoteric contexts like algebraic geometry over

finite fields which at first sight lack the differential structure needed for constructions such

as tangent vector.

**Another useful application of locally ringed spaces is in the construction of schemes. The
**

forgetful functor assigning to each locally ringed space (X, OX ) the ring OX (X) is adjoint to

the ”prime spectrum” functor taking each ring R to its prime spectrum Spec(R), and this

correspondence is essentially why the category of locally ringed spaces is the proper building

block to use in the formulation of the notion of scheme.

Version: 9 Owner: djao Author(s): djao

292.2 presheaf

**For a topological space X a presheaf F with values in a category C associates to each open set
**

U ⊂ X, an object F (U) of C and to each inclusion U ⊂ V a morphism of C, ρU V : F (V ) →

F (U), the restriction morphism. It is required that ρU U = 1F (U ) and ρU W = ρU V ◦ ρV W for

any U ⊂ V ⊂ W .

**A presheaf with values in the category of sets (or abelian groups) is called a presheaf of sets
**

(or abelian groups). If no target category is specified, either the category of sets or abelian

groups is most likely understood.

**A more categorical way to state it is as follows. For X form the category Top(X) whose
**

objects are open sets of X and whose morphisms are the inclusions. Then a presheaf is

merely a contravariant functor Top(X) → C.

Version: 2 Owner: nerdy2 Author(s): nerdy2

292.3 sheaf

292.3.1 Presheaves

**Let X be a topological space and let A be a category. A presheaf on X with values in A is
**

a contravariant functor F from the category of open sets in X and inclusion morphisms to

the category A.

As this definition may be less than helpful to many readers, we offer the following equivalent

1220

(but longer) definition. A presheaf F on X consists of the following data:

1. An object F (U) in A, for each open set U ⊂ X

**2. A morphism resV,U : F (V ) −→ F (U) for each pair of open sets U ⊂ V in X (called
**

the restriction morphism), such that:

**(a) For every open set U ⊂ X, the morphism resU,U is the identity morphism.
**

(b) For any open sets U ⊂ V ⊂ W in X, the diagram

resW,U

F (W ) resW,V F (V ) resV,U F (U)

commutes.

If the object F (U) of A is a set, its elements are called sections of U.

292.3.2 Morphisms of Presheaves

**Let f : X −→ Y be a continuous map of topological spaces. Suppose FX is a presheaf on
**

X, and GY is a presheaf on Y (with FX and GY both having values in A). We define a

morphism of presheaves φ from GY to FX , relative to f , to be a collection of morphisms

φU : GY (U) −→ FX (f −1 (U)) in A, one for every open set U ⊂ Y , such that the diagram

φV

GY (V ) FX (f −1 (V ))

resV,U resf −1 (V ),f −1 (U )

GY (U) φU

FX (f −1 (U))

commutes, for each pair of open sets U ⊂ V in Y .

**In the special case that f is the identity map id : X −→ X, we omit mention of the map f ,
**

and speak of φ as simply a morphism of presheaves on X. Form the category whose objects

are presheaves on X and whose morphisms are morphisms of presheaves on X. Then an

isomorphism of presheaves φ on X is a morphism of presheaves on X which is an isomorphism

in this category; that is, there exists a morphism φ−1 whose composition with φ both ways

is the identity morphism.

**More generally, if f : X −→ Y is any homeomorphism of topological spaces, a morphism of
**

presheaves φ relative to f is an isomorphism if it admits a two–sided inverse morphism of

presheaves φ−1 relative to f −1 .

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292.3.3 Sheaves

**We now assume that the category A is a concrete category. A sheaf is a presheaf F on X,
**

with values in A, such that for every open set U ⊂ X, and every open cover {Ui } of U, the

following two conditions hold:

**1. Any two elements f1 , f2 ∈ F (U) which have identical restrictions to each Ui are equal.
**

That is, if resU,Ui f1 = resU,Ui f2 for every i, then f1 = f2 .

**2. Any collection of elements fi ∈ F (Ui ) that have common restrictions can be realized
**

as the collective restrictions of a single element of F (U). That is, if resUi ,Ui T Uj fi =

resUj ,Ui T Uj fj for every i and j, then there exists an element f ∈ F (U) such that

resU,Ui f = fi for all i.

292.3.4 Sheaves in abelian categories

**If A is a concrete abelian category, then a presheaf F is a sheaf if and only if for every open
**

subset U of X, the sequence

incl Q diff

Q T

0 F (U) i F (Ui ) i,j F (Ui Uj ) (292.3.1)

**is an exact sequence of morphisms in A for every open cover {Ui } of U in X. This diagram
**

requires some explanation, because we owe the reader a definition of the morphisms incl and

diff. We start with incl (short for “inclusion”). The restriction morphisms F (U) −→ F (Ui )

induce a morphism Y

F (U) −→ F (Ui )

i

Q

to the categorical direct product i F (Ui ), which we define to be incl. The map diff (called

“difference”) is defined as follows. For each Ui , form the morphism

Y \

αi : F (Ui ) −→ F (Ui Uj ).

j

**By the universal properties of categorical direct product, there exists a unique morphism
**

Y YY \

α: F (Ui ) −→ F (Ui Uj )

i i j

**such that πi α = αi πi for all i, where πi is projection onto the ith factor. In a similar manner,
**

form the morphism Y YY \

β: F (Uj ) −→ F (Ui Uj ).

j j i

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Then α and β are both elements of the set

!

Y Y \

Hom F (Ui ), F (Ui Uj ) ,

i i,j

**which is an abelian group since A is an abelian category. Take the difference α − β in this
**

group, and define this morphism to be diff.

**Note that exactness of the sequence (291.3.1) is an element free condition, and therefore
**

makes sense for any abelian category A, even if A is not concrete. Accordingly, for any

abelian category A, we define a sheaf to be a presheaf F for which the sequence (291.3.1) is

always exact.

292.3.5 Examples

**It’s high time that we give some examples of sheaves and presheaves. We begin with some
**

of the standard ones.

Example 9. If F is a presheaf on X, and U ⊂ X is an open subset, then one can define a

presheaf F |U on U by restricting the functor F to the subcategory of open sets of X in U

and inclusion morphisms. In other words, for open subsets of U, define F |U to be exactly

what F was, and ignore open subsets of X that are not open subsets of U. The resulting

presheaf is called, for obvious reasons, the restriction presheaf of F to U, or the restriction

sheaf if F was a sheaf to begin with.

Example 10. For any topological space X, let cX be the presheaf on X, with values in the

category of rings, given by

**• cX (U) := the ring of continuous real–valued functions U −→ R,
**

• resV,U f := the restriction of f to U, for every element f : V −→ R of cX (V ) and every

subset U of V .

**Then cX is actually a sheaf of rings, because continuous functions are uniquely specified by
**

their values on an open cover. The sheaf cX is called the sheaf of continuous real–valued

functions on X.

Example 11. Let X be a smooth differentiable manifold. Let DX be the presheaf on X, with

values in the category of real vector spaces, defined by setting DX (U) to be the space of

smooth real–valued functions on U, for each open set U, and with the restriction morphism

given by restriction of functions as before. Then DX is a sheaf as well, called the sheaf of

smooth real–valued functions on X.

**Much more surprising is that the construct DX can actually be used to define the concept
**

of smooth manifold! That is, one can define a smooth manifold to be a locally Euclidean

n–dimensional second countable topological space X, together with a sheaf F , such that

there exists an open cover {Ui } of X where:

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For every i, there exists a homeomorphism fi : Ui −→ Rn and an isomorphism

of sheaves φi : DRn −→ F |Ui relative to fi .

The idea here is that not only does every smooth manifold X have a sheaf DX of smooth functions,

but specifying this sheaf of smooth functions is sufficient to fully describe the smooth man-

ifold structure on X. While this phenomenon may seem little more than a toy curiousity

for differential geometry, it arises in full force in the field of algebraic geometry where the

coordinate functions are often unwieldy and algebraic structures in many cases can only be

satisfactorily described by way of sheaves and schemes.

Example 12. Similarly, for a complex analytic manifold X, one can form the sheaf HX of

holomorphic functions by setting HX (U) equal to the complex vector space of C–valued

holomorphic functions on U, with the restriction morphism being restriction of functions as

before.

Example 13. The algebraic geometry analogue of the sheaf DX of differential geometry is the

prime spectrum Spec(R) of a commutative ring R. However, the construction of the sheaf

Spec(R) is beyond the scope of this discussion and merits a separate article.

Example 14. For an example of a presheaf that is not a sheaf, consider the presheaf F on X,

with values in the category of real vector spaces, whose sections on U are locally constant

real–valued functions on U modulo constant functions on U. Then every section f ∈ F (U)

is locally zero in some fine enough open cover {Ui } (it is enough to take a cover where each

Ui is connected), whereas f may be nonzero if U is not connected.

**We conclude with some interesting examples of morphisms of sheaves, chosen to illustrate the
**

unifying power of the language of schemes across various diverse branches of mathematics.

**1. For any continuous function f : X −→ Y , the map φU : cY (U) −→ cX (f −1 (U)) given
**

by φU (g) := gf defines a morphisms of sheaves from cY to cX with respect to f .

2. For any continuous function f : X −→ Y of smooth differentiable manifolds, the map

given by φU (g) := gf has the property

g ∈ DY (U) → φU (g) ∈ DX (f −1 (U))

if and only if f is a smooth function.

3. For any continuous function f : X −→ Y of complex analytic manifolds, the map given

by φU (g) := gf has the property

g ∈ HY (U) → φU (g) ∈ HX (f −1 (U))

if and only if f is a holomorphic function.

4. For any Zariski continuous function f : X −→ Y of algebraic varieties over a field k,

the map given by φU (g) := gf has the property

g ∈ OY (U) → φU (g) ∈ OX (f −1 (U))

if and only if f is a regular function. Here OX denotes the sheaf of k–valued regular

functions on the algebraic variety X.

1224

REFERENCES

1. David Mumford, The Red Book of Varieties and Schemes, Second Expanded Edition, Springer–

Verlag, 1999 (LNM 1358).

2. Charles Weibel, An Introduction to Homological Algebra, Cambridge University Press, 1994.

Version: 9 Owner: djao Author(s): djao

292.4 sheafification

**Let F be a presheaf over a topological space X with values in a category A for which sheaves
**

are defined. The sheafification of F , if it exists, is a sheaf F 0 over X together with a morphism

θ : F −→ F 0 satisfying the following universal property:

**For any sheaf G over X and any morphism of presheaves φ : F −→ G over X,
**

there exists a unique morphism of sheaves ψ : F 0 −→ G such that the diagram

θ ψ

F F0 G

φ

commutes.

**In light of the universal property, the sheafification of F is uniquely defined up to canonical
**

isomorphism whenever it exists. In the case where A is a concrete category (one consisting

of sets and set functions), the sheafification ofSany presheaf F can be constructed by taking

F 0 (U) to be the set of all functions s : U −→ p∈U Fp such that

1. s(p) ∈ Fp for all p ∈ U

**2. For all p ∈ U, there is a neighborhood V ⊂ U of p and a section t ∈ F (V ) such that,
**

for all q ∈ V , the induced element tq ∈ Fq equals s(q)

for all open sets U ⊂ X. Here Fp denotes the stalk of the presheaf F at the point p.

**The following quote, taken from [1], is perhaps the best explanation of sheafification to be
**

found anywhere:

**F 0 is ”the best possible sheaf you can get from F ”. It is easy to imagine how
**

to get it: first identify things which have the same restrictions, and then add in

all the things which can be patched together.

1225

REFERENCES

1. David Mumford, The Red Book of Varieties and Schemes, Second Expanded Edition, Springer–

Verlag, 1999 (LNM 1358)

Version: 4 Owner: djao Author(s): djao

292.5 stalk

**Let F be a presheaf over a topological space X with values in an abelian category A, and
**

suppose direct limits exist in A. For any point p ∈ X, the stalk Fp of F at p is defined to

be the object in A which is the direct limit of the objects F (U) over the directed set of all

open sets U ⊂ X containing p, with respect to the restriction morphisms of F . In other

words,

Fp := lim F (U)

−→

U 3p

**If A is a category consisting of sets, the stalk Fp can be viewed as the set of all germs of
**

sections of F at the point p. That is, the set Fp consists of all the equivalence classes of

ordered pairs (U, s) where p ∈ U and s T∈ F (U), under the equivalence relation (U, s) ∼ (V, t)

if there exists a neighborhood W ⊂ U V of p such that resU,W s = resV,W t.

**By universal properties of direct limit, a morphism φ : F −→ G of presheaves over X induces
**

a morphism φp : Fp −→ Gp on each stalk Fp of F . Stalks are most useful in the context

of sheaves, since they encapsulate all of the local data of the sheaf at the point p (recall

that sheaves are basically defined as presheaves which have the property of being completely

characterized by their local behavior). Indeed, in many of the standard examples of sheaves

that take values in rings (such as the sheaf DX of smooth functions, or the sheaf OX of

regular functions), the ring Fp is a local ring, and much of geometry is devoted to the study

of sheaves whose stalks are local rings (so-called ”locally ringed spaces”).

**We mention here a few illustrations of how stalks accurately reflect the local behavior of a
**

sheaf; all of these are drawn from [1].

**• A morphism of sheaves φ : F −→ G over X is an isomorphism if and only if the
**

induced morphism φp is an isomorphism on each stalk.

**• A sequence F −→ G −→ H of morphisms of sheaves over X is an exact sequence at
**

G if and only if the induced morphism Fp −→ Gp −→ Hp is exact at each stalk Gp .

• The sheafification F 0 of a presheaf F has stalk equal to Fp at every point p.

1226

REFERENCES

1. Robin Hartshorne, Algebraic Geometry, Springer–Verlag New York Inc., 1977 (GTM 52).

Version: 4 Owner: djao Author(s): djao

1227

Chapter 293

18F30 – Grothendieck groups

293.1 Grothendieck group

**Let S be an Abelian semigroup. The Grothendieck group of S is K(S) = S × S/∼, where
**

∼ is the equivalence relation: (s, t) ∼ (u, v) if there exists r ∈ S such that s + v + r =

t + u + r. This is indeed an abelian group with zero element (s, s) (any s ∈ S) and inverse

−(s, t) = (t, s).

**The Grothendieck group construction is a functor from the category of abelian semigroups
**

to the category of abelian groups. A morphism f : S → T induces a morphism K(f ) :

K(S) → K(T ).

Example 15. K(N) = Z.

Example 16. Let G be an abelian group, then K(G) ∼

= G via (g, h) ↔ g − h.

**Let C be a symmetric monoidal category. Its Grothendieck group is K([C]), i.e. the
**

Grothendieck group of the isomorphism classes of objects of C.

Version: 2 Owner: mhale Author(s): mhale

1228

Chapter 294

18G10 – Resolutions; derived functors

294.1 derived functor

There are two objects called derived functors. First, there are classical derived functors. Let

A, B be abelian categories, and F : A → B be a covariant left-exact functor. Note that a

completely analogous construction can be done for right-exact and contravariant functors,

but it is traditional to only describe one case, as doing the other mostly consists of reversing

arrows. Given an object A ∈ A, we can construct an injective resolution:

A→A:0 A I1 I2 ···

**which is unique up to chain homotopy equivalence. Then we apply the functor F to the
**

injectives in the resolution to to get a complex

F (A) : 0 F (I 1 ) F (I 2 ) ···

**(notice that the term involving A has been left out. This is not an accident, in fact, it is
**

crucial). This complex also is independent of the choice of I’s (up to chain homotopy equiv-

alence). Now, we define the classical right derived functors Ri F (A) to be the cohomology

groups H i (F (A)). These only depend on A.

**Important properties of the classical derived functors are these: If the sequence 0 → A →
**

A0 → A00 → 0 is exact, then there is a long exact sequence

0 F (A) F (A0 ) F (A00 ) R1 F (A) ···

**which is natural (a morphism of short exact sequences induces a morphism of long exact
**

sequences). This, along with a couple of other properties determine the derived functors

completely, giving an axiomatic definition, though the construction used above is usually

necessary to show existence.

From the definition, one can see immediately that the following are equivalent:

1229

1. F is exact

2. Rn F (A) = 0 for n > 1 and all A ∈ A.

3. R1 F (A) = 0 for all A ∈ A.

However, R1 F (A) = 0 for a particular A does not imply that Rn F (A) = 0 for all n > 1.

**Important examples are Extn , the derived functor of Hom, Torn , the derived functor of
**

tensor product, and sheaf cohomology, the derived functor of the global section functor on

sheaves.

(Coming soon: the derived categoies definition)

Version: 4 Owner: bwebste Author(s): bwebste

1230

Chapter 295

**18G15 – Ext and Tor, generalizations,
**

Künneth formula

295.1 Ext

**For a ring R, and R-module A, we have a covariant functor HomA − R. ExtnR (A, −) are
**

defined to be the right derived functors of HomA − R (ExtnR (A, −) = Rn HomA − R).

**Ext gets its name from the following fact: There is a natural bijection between elements of
**

Ext1R (A, B) and extensions of B by A up to isomorphism of short exact sequences, where an

extension of B by A is an exact sequence

0→B→C→A→0

. For example,

Ext1Z (Z/nZ, Z) ∼

= Z/nZ

, with 0 corresponding to the trivial extension 0 → Z → Z ⊕ Z/nZ → 0, and m 6= 0

corresponding to

n m

0 Z Z Z/nZ 0.

Version: 3 Owner: bwebste Author(s): bwebste

1231

Chapter 296

**18G30 – Simplicial sets, simplicial
**

objects (in a category)

296.1 nerve

The

**Definition 18. nerve of a category C is the simplicial set Hom(i(−), C), where i : ∆ → Cat
**

is the fully faithful functor that takes each ordered set [n] in the simplicial category, ∆, to

op

the pre-order n + 1. The nerve is a functor Cat → Set∆ .

Version: 1 Owner: mhale Author(s): mhale

296.2 simplicial category

**The simplicial category ∆ is defined as the small category whose objects are the totally ordered
**

finite sets

[n] = {0 < 1 < 2 < . . . < n}, n > 0, (296.2.1)

and whose morphisms are monotonic non-decreasing (order-preserving) maps. It is generated

by two families of morphisms:

**δin : [n − 1] → [n] is the injection missing i ∈ [n],
**

σin : [n + 1] → [n] is the surjection such that σin (i) = σin (i + 1) = i ∈ [n].

The δin morphisms are called

Definition 19. face maps, and the σin morphisms are called

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Definition 20. degeneracy maps. They satisfy the following relations,

**δjn+1 δin = δin+1 δj−1
**

n

for i < j, (296.2.2)

σjn−1 σin = σin−1 σj+1

n

for i 6 j, (296.2.3)

n n−1

δi σj−1 if i < j,

n n+1

σj δi = idn if i = j or i = j + 1, (296.2.4)

n n−1

δi−1 σj if i > j + 1.

All morphisms [n] → [0] factor through σ00 , so [0] is terminal.

There is a bifunctor + : ∆ × ∆ → ∆ defined by

[m] + [n] = [m + n + 1], (296.2.5)

f (i) if 0 6 i 6 m,

(f + g)(i) = (296.2.6)

g(i − m − 1) + m0 + 1 if m < i 6 (m + n + 1),

**where f : [m] → [m0 ] and g : [n] → [n0 ]. Sometimes, the simplicial category is defined to
**

include the empty set [−1] = ∅, which provides an initial object for the category. This makes

∆ a strict monoidal category as ∅ is a unit for the bifunctor: ∅ + [n] = [n] = [n] + ∅ and

id∅ + f = f = f + id∅ . Further, ∆ is then the free monoidal category on a monoid object

(the monoid object being [0], with product σ00 : [0] + [0] → [0]).

**There is a fully faithful functor from ∆ to Top, which sends each object [n] to an oriented n-
**

simplex. The face maps then embed an (n − 1)-simplex in an n-simplex, and the degeneracy

maps collapse an (n + 1)-simplex to an n-simplex. The bifunctor forms a simplex from the

disjoint union of two simplicies by joining their vertices together in a way compatible with

their orientations.

**There is also a fully faithful functor from ∆ to Cat, which sends each object [n] to a pre-order
**

n + 1. The pre-order n is the category consisting of n partially-ordered objects, with one

morphism a → b iff a 6 b.

Version: 4 Owner: mhale Author(s): mhale

296.3 simplicial object

A

**Definition 21. simplicial object in a category C is a contravariant functor from the simplicial category
**

∆ to C. Such a functor X is uniquely specified by the morphisms X(δin ) : [n] → [n − 1] and

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X(σin ) : [n] → [n + 1], which satisfy

**X(δin−1 ) X(δjn ) = X(δj−1
**

n−1

) X(δin ) for i < j, (296.3.1)

X(σin+1 ) X(σjn ) = X(σj+1

n+1

) X(σin ) for i 6 j, (296.3.2)

n−1

X(σj−1 ) X(δin ) if i < j,

X(δin+1 ) X(σjn ) = idn if i = j or i = j + 1, (296.3.3)

X(σjn−1 ) X(δi−1

n

) if i > j + 1.

In particular, a

**Definition 22. simplicial set is a simplicial object in Set. Equivalently, one could say that a
**

simplicial set is a presheaf on ∆. The object X([n]) of a simplicial set is a set of n-simplices,

and is called the n-skeleton.

Version: 2 Owner: mhale Author(s): mhale

1234

Chapter 297

18G35 – Chain complexes

297.1 5-lemma

**If Ai , Bi for i = 1, . . . , 5 are objects in an abelian category (for example, modules over a ring
**

R) such that there is a commutative diagram

A1 A2 A3 A4 A5

γ1 γ2 γ3 γ4 γ5

B1 B2 B3 B4 B5

**with the rows exact, and γ1 is surjective, γ5 is injective, and γ2 and γ4 are isomorphisms,
**

then γ3 is an isomorphism as well.

Version: 2 Owner: bwebste Author(s): bwebste

1235

297.2 9-lemma

**If Ai , Bi , Ci , for i = 1, 2, 3 are objects of an abelian category such that there is a commutative diagram
**

0 0 0

0 A1 B1 C1 0

0 A2 B2 C2 0

0 A3 B3 C3 0

0 0 0

with the columns and bottom two rows are exact, then the top row is exact as well.

Version: 2 Owner: bwebste Author(s): bwebste

297.3 Snake lemma

**There are two versions of the snake lemma: (1) Given a commutative (1) diagram as below,
**

with exact (1) rows

0 −−−→ A1 −−−→ B1 −−−→ C1 −−−→ 0

αy

βy

γ y

0 −−−→ A2 −−−→ B2 −−−→ C2 −−−→ 0

there is an exact sequence

0 → ker α → ker β → ker γ → coker α → coker β → coker γ → 0

where ker denotes the kernel of a map and coker its cokernel.

**(2) Applying this result inductively to a short exact (2) sequence of (2) chain complexes, we
**

obtain the following: Let A, B, C be chain complexes, and let

0→A→B→C→0

be a short exact sequence. Then there is a long exact sequence of homology groups

· · · → Hn (A) → Hn (B) → Hn (C) → Hn−1 (A) → · · ·

Version: 5 Owner: bwebste Author(s): bwebste

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297.4 chain homotopy

0 0 0 0

Let (A, d) and (A , d ) be chain complexes and f : A → A , g : A → A be chain maps. A

0

chain homotopy D between f and g is a sequence of homomorphisms {Dn : An → An+1 }

0

so that dn+1 ◦ Dn + Dn−1 ◦ dn = fn − gn for each n. Thus, we have a commutative diagram:

dn+1 dn

An+1 An An−1

Dn−1

fn+1 −gn+1 fn−1 −gn−1

Dn

0 0 0

An+1 0 An 0

An−1

dn+1 dn

Version: 4 Owner: RevBobo Author(s): RevBobo

**297.5 chain map
**

0 0 0

Let (A, d) and (A , d ) be chain complexes. A chain map f : A → A is a sequence of

0

homomorphisms {fn } such that dn ◦ fn = fn−1 ◦ dn for each n. Diagramatically, this says

that the following diagram commutes:

dn

An An−1

fn fn−1

0

0 dn 0

An An−1

Version: 3 Owner: RevBobo Author(s): RevBobo

297.6 homology (chain complex)

**If (A, d) is a chain complex
**

dn−1 d dn+1 dn+2

· · · ←−−− An−1 ←−n−− An ←−−− An+1 ←−−− · · ·

then the n-th homology group Hn (A, d) (or module) of the chain complex A is the quotient

ker dn

Hn (A, d) = .

i dn+1

Version: 2 Owner: bwebste Author(s): bwebste

1237

Chapter 298

**18G40 – Spectral sequences,
**

hypercohomology

298.1 spectral sequence

**A spectral sequence is a collection of R-modules (or more generally, objects of an abelian category)
**

r

{Ep,q r

} for all r ∈ N, p, q ∈ Z, equipped with maps drpq : Ep,q r

→ Ep−r,q+r−1 such that is a

r+1

chain complex, and the E ’s are its homology, that is,

r+1 ∼

Ep,q = ker(drp,q )/im(drp+r,q−r+1).

**(Note: what I have defined above is a homology spectral sequence. Cohomology spectral
**

sequences are identical, except that all the arrows go in the other direction.)

r

Most interesting spectral sequences are upper right quadrant, meaning that Ep,q = 0 if p

or q < 0. If this is the case then for any p, q, both drpq and drp+r,q−r+1 are 0 for sufficiently

large r since the target or source is out of the upper right quadrant, so that for all r > r0

r r+1 ∞

Ep,q = Ep,q · · · . This group is called Ep,q .

r

A upper right quadrant spectral sequence {Ep,q } is said to converge to a sequence Fn of

R-modules if there is an exhaustive filtration Fn,0 = 0 ⊂ Fn,1 ⊂ · · · ⊂ of each Fn such that

Fp+q,q+1 /Fp+q,q ∼ ∞

= Ep,q

r

. This is typically written Ep,q ⇒ Fp+q .

**Typically spectral sequences are used in the following manner: we find an interpretation of
**

E r for a small value of r, typically 1, and of E ∞ , and then in cases where enough groups

and differentials are 0, we can obtain information about one from the other.

Version: 2 Owner: bwebste Author(s): bwebste

1238

Chapter 299

**19-00 – General reference works
**

(handbooks, dictionaries,

bibliographies, etc.)

299.1 Algebraic K-theory

**Algebraic K-theory is a series of functors on the category of rings. It classifies ring invariants,
**

i.e. ring properties that are Morita invariant.

The functor K0

**Let R be a ring and denote by M∞ (R) the algebraic direct
**

limit

of matrix algebras Mn (R)

a 0

under the embeddings Mn (R) → Mn+1 (R) : a → . The zeroth K-group of

0 0

R, K0 (R), is the Grothendieck group (abelian group of formal differences) of the unitary

equivalence classes of projections in M∞ (R). The addition of two equivalence classes [p] and

[q] is given by the direct summation of the projections p and q: [p] + [q] = [p ⊕ q].

The functor K1

[To Do: coauthor?]

The functor K2

[To Do: coauthor?]

Higher K-functors

Higher K-groups are defined using the Quillen plus construction,

Knalg (R) = πn (BGL∞ (R)+ ), (299.1.1)

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where GL∞ (R) is the infinite general linear group over R (defined in a similar way to

M∞ (R)), and BGL∞ (R) is its classifying space.

Algebraic K-theory has a product structure,

Ki (R) ⊗ Kj (S) → Ki+j (R ⊗ S). (299.1.2)

Version: 2 Owner: mhale Author(s): mhale

299.2 K-theory

**Topological K-theory is a generalised cohomology theory on the category of compact Hausdorff
**

spaces. It classifies the vector bundles over a space X up to stable equivalences. Equivalently,

via the Serre-Swan theorem, it classifies the finitely generated projective modules over the

C ∗ -algebra C(X).

**Let A be a unital C ∗ -algebra over C and denote by M∞ (A) the algebraicdirect limit of
**

a 0

matrix algebras Mn (A) under the embeddings Mn (A) → Mn+1 (A) : a → . The

0 0

K0 (A) group is the Grothendieck group (abelian group of formal differences) of the homotopy

classes of the projections in M∞ (A). Two projections p and q are homotopic if p = uqu−1 for

some unitary u ∈ M∞ (A). Addition of homotopy classes is given by the direct summation

of projections: [p] + [q] = [p ⊕ q].

**Denote by U∞ (A) the direct
**

limit of unitary groups Un (A) under the embeddings Un (A) →

u 0

Un+1 (A) : u → . Give U∞ (A) the direct limit topology, i.e. a subset U of U∞ (A)

0 1 T

is open if and only if U Un (A) is an open subset of Un (A), for all n. The K1 (A) group

is the Grothendieck group (abelian group of formal differences) of the homotopy classes of

the unitaries in U∞ (A). Two unitaries u and v are homotopic if uv −1 lies in the identity

component of U∞ (A). Addition of homotopy classes is given by the direct summation of

unitaries: [u] + [v] = [u ⊕ v]. Equivalently, one can work with invertibles in GL∞ (A) (an

invertible g is connected to the unitary u = g|g|−1 via the homotopy t → g|g|−t).

Higher K-groups can be defined through repeated suspensions,

Kn (A) = K0 (S n A). (299.2.1)

But, the Bott periodicity theorem means that

K1 (SA) ∼

= K0 (A). (299.2.2)

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The main properties of Ki are:

Ki (A ⊕ B) = Ki (A) ⊕ Ki (B), (299.2.3)

Ki (Mn (A)) = Ki (A) (Morita invariance), (299.2.4)

Ki (A ⊗ K) = Ki (A) (stability), (299.2.5)

Ki+2 (A) = Ki (A) (Bott periodicity). (299.2.6)

**There are three flavours of topological K-theory to handle the cases of A being complex (over
**

C), real (over R) or Real (with a given real structure).

**Ki (C(X, C)) = KU −i (X) (complex/unitary), (299.2.7)
**

Ki (C(X, R)) = KO −i (X) (real/orthogonal), (299.2.8)

KR i (C(X), J) = KR −i (X, J) (Real). (299.2.9)

Real K-theory has a Bott period of 8, rather than 2.

REFERENCES

1. N. E. Wegge-Olsen, K-theory and C ∗ -algebras. Oxford science publications. Oxford University

Press, 1993.

2. B. Blackadar, K-Theory for Operator Algebras. Cambridge University Press, 2nd ed., 1998.

Version: 12 Owner: mhale Author(s): mhale

299.3 examples of algebraic K-theory groups

R K0 (R) K1 (R) K2 (R) K3 (R) K4 (R)

Z Z Z/2 Z/2 Z/48 0

R Z R×

C C C×

Algebraic K-theory of some common rings.

Version: 2 Owner: mhale Author(s): mhale

1241

Chapter 300

19K33 – EXT and K-homology

300.1 Fredholm module

**Fredholm modules represent abstract elliptic pseudo-differential operators.
**

Definition 3. An

Definition 23. odd Fredholm module (H, F ) over a C ∗ -algebra A is given by an involutive

representation π of A on a Hilbert space H, together with an operator F on H such that

F = F ∗ , F 2 = 1I and [F, π(a)] ∈ K(H) for all a ∈ A.

Definition 4. An

Definition 24. even Fredholm module (H, F, Γ) is given by an odd Fredholm module (H, F )

together with a Z2 -grading Γ on H, Γ = Γ∗ , Γ2 = 1I, such that Γπ(a) = π(a)Γ and ΓF =

−F Γ.

Definition 5. A Fredholm module is called

Definition 25. degenerate if [F, π(a)] = 0 for all a ∈ A. Degenerate Fredholm modules are

homotopic to the 0-module.

Example 17 (Fredholm modules over C). An even Fredholm module (H, F, Γ) over C is given

by

k k a1Ik 0

H = C ⊕ C with π(a) = ,

0 0

0 1Ik

F = ,

1Ik 0

1Ik 0

Γ = .

0 −1Ik

Version: 3 Owner: mhale Author(s): mhale

1242

300.2 K-homology

**K-homology is a homology theory on the category of compact Hausdorff spaces. It classifies
**

the elliptic pseudo-differential operators acting on the vector bundles over a space. In terms

of C ∗ -algebras, it classifies the Fredholm modules over an algebra.

**The K 0 (A) group is the abelian group of homotopy classes of even Fredholm modules over
**

A. The K 1 (A) group is the abelian group of homotopy classes of odd Fredholm modules over

A. Addition is given by direct summation of Fredholm modules, and the inverse of (H, F, Γ)

is (H, −F, −Γ).

Version: 1 Owner: mhale Author(s): mhale

1243

Chapter 301

19K99 – Miscellaneous

301.1 examples of K-theory groups

A K0 (A) K1 (A)

C Z 0

Mn (C) Z 0

H Z 0

K Z 0

B 0 0

B/K 0 Z

C0 ((0, 1)) 0 Z

C0 (R2n ) Z 0

C0 (R2n+1 ) 0 Z

C([0, 1]) Z 0

n−1 n−1

C(Tn ) Z2 Z2

C(S2n ) Z2 0

C(S2n+1 ) Z Z

C(CPn ) Zn+1 0

On Z/(n − 1) 0

Aθ Z2 Z2

C ∗ (H3 ) Z3 Z3

Topological K-theory of some common C ∗ -algebras.

Version: 5 Owner: mhale Author(s): mhale

1244

Chapter 302

**20-00 – General reference works
**

(handbooks, dictionaries,

bibliographies, etc.)

**302.1 alternating group is a normal subgroup of the
**

symmetric group

Theorem 2. The alternating group An is a normal subgroup of the symmetric group Sn

**D efine the epimorphism f : Sn → Z2 by : σ 7→ 0 if σ is an even permutation and : σ 7→ 1
**

if σ is an odd permutation. Hence, An is the kernel of f and so it is a normal subgroup

of the domain Sn . Furthermore Sn /An ∼ = Z2 by the first isomorphism theorem. So by

Lagrange’s theorem

|Sn | = |An ||Sn /An |.

Therefore, |An | = n!/2. That is, there are n!/2 many elements in An

Version: 1 Owner: tensorking Author(s): tensorking

302.2 associative

Let (S, φ) be a set with binary operation φ. φ is said to be associative over S if

φ(a, φ(b, c)) = φ(φ(a, b), c)

1245

for all a, b, c ∈ S.

**Examples of associative operations are addition and multiplication over the integers (or
**

reals), or addition or multiplication over n × n matrices.

**We can construct an operation which is not associative. Let S be the integers. and define
**

ν(a, b) = a2 + b. Then ν(ν(a, b), c) = ν(a2 + b, c) = a4 + 2ba2 + b2 + c. But ν(a, ν(b, c)) =

ν(a, b2 + c) = a + b4 + 2cb2 + c2 , hence ν(ν(a, b), c) 6= ν(a, ν(b, c)).

**Note, however, that if we were to take S = {0}, ν would be associative over S!. This
**

illustrates the fact that the set the operation is taken with respect to is very important.

**Example. We show that the division operation over nonzero reals is non-associative. All
**

we need is a counter-example: so let us compare 1/(1/2) and (1/1)/2. The first expression

is equal to 2, the second to 1/2, hence division over the nonzero reals is not associative.

Version: 6 Owner: akrowne Author(s): akrowne

302.3 canonical projection

Given a group G and a normal subgroup N G there is an epimorphism

π : G → G/N

**defined by sending an element g ∈ G to its coset gN. The epimorphism π is referred to as
**

the canonical projection.

Version: 4 Owner: Dr Absentius Author(s): Dr Absentius

302.4 centralizer

For a given group G, the centralizer of an element a ∈ G is defined to be the set

C(a) = {x ∈ G | xa = ax}

**We note that if x, y ∈ C(a) then xy −1 a = xay −1 = axy −1 so that xy −1 ∈ C(a). Thus C(a)
**

is a non-trivial subgroup of G containing at least {e, a}. To illustrate an application of this

concept we prove the following lemma.

There exists a bijection between the right cosets of C(a) and the conjugates of a.

**If x, y ∈ G are in the same right coset, then y = cx for some c ∈ C(a). Thus y −1ay =
**

x−1 c−1 acx = x−1 c−1 cax = x−1 ax. Conversely, if y −1 ay = x−1 ax then xy −1 a = axy −1 and

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xy −1 ∈ C(a) giving x, y are in the same right coset. Let [a] denote the conjugacy class of a.

It follows that |[a]| = [G : C(a)] and |[a]| | |G|.

We remark that a ∈ Z(G) ⇔ |C(a)| = |G| ⇔ |[a]| = 1, where Z(G) denotes the center of G.

**Now let G be a p-group, i.e. a finite group of order pn , where p is a prime and n P > 0. Let
**

z =P |Z(G)|. Summing over elements in distinct conjugacy classes, we have pn = |[a]| =

z + a∈Z(G)

/ |[a]| since the center consists precisely of the conjugacy classes of cardinality 1.

n

But |[a]| | p , so p | z. However, Z(G) is certainly non-empty, so we conclude that every

p-group has a non-trivial center.

The groups C(gag −1) and C(a), for any g, are isomorphic.

Version: 5 Owner: mathcam Author(s): Larry Hammick, vitriol

302.5 commutative

Let (S, φ) be a set with binary operation φ. φ is said to be commutative if

φ(a, b) = φ(b, a)

for all a, b ∈ S.

**Some operations which are commutative are addition over the integers, multiplication over
**

the integers, addition over n × n matrices, and multiplication over the reals.

An example of a non-commutative operation is multiplication over n × n matrices.

Version: 3 Owner: akrowne Author(s): akrowne

302.6 examples of groups

**Groups are ubiquitous throughout mathematics. Many “naturally occurring” groups are
**

either groups of numbers (typically abelian) or groups of symmetries (typically non-abelian).

Groups of numbers

**• The most important group is the group of integers Z with addition as operation.
**

• The integers modulo n, often denoted by Zn , form a group under addition. Like Z

itself, this a cyclic group; any cyclic group is isomorphic to one of these.

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• The rational (or real, or complex) numbers form a group under addition.

**• The positive rationals form a group under multiplication, and so do the non-zero
**

rationals. The same is true for the reals.

**• The non-zero complex numbers form a group under multiplication. So do the non-zero
**

quaternions. The latter is our first example of a non-abelian group.

• More generally, any (skew) field gives rise to two groups: the additive group of all field

elements, and the multiplicative group of all non-zero field elements.

**• The complex numbers of absolute value 1 form a group under multiplication, best
**

thought of as the unit circle. The quaternions of absolute value 1 form a group under

multiplication, best thought of as the three-dimensional unit sphere S 3 . The two-

dimensional sphere S 2 however is not a group in any natural way.

Most groups of numbers carry natural topologies turning them into topological groups.

Symmetry groups

**• The symmetric group of degree n, denoted by Sn , consists of all permutations of n
**

items and has n! elements. Every finite group is isomorphic to a subgroup of some Sn .

**• An important subgroup of the symmetric group of degree n is the alternating group,
**

denoted An . This consists of all even permutations on n items. A permutation is said

to be even if it can be written as the product of an even number of transpositions.

The alternating group is normal in Sn , of index 2, and it is an interesting fact that An

is simple for n > 5. See the proof on the simplicity of the alternating groups. By the

Jordan-Hölder theorem, this means that this is the only normal subgroup of Sn .

**• If any geometrical object is given, one can consider its symmetry group consisting
**

of all rotations and reflections which leave the object unchanged. For example, the

symmetry group of a cone is isomorphic to S 1 .

**• The set of all automorphisms of a given group (or field, or graph, or topological space,
**

or object in any category) forms a group with operation given by the composition of

homomorphisms. These are called automorphism groups; they capture the internal

symmetries of the given objects.

**• In Galois theory, the symmetry groups of field extensions (or equivalently: the symme-
**

try groups of solutions to polynomial equations) are the central object of study; they

are called Galois groups.

• Several matrix groups describe various aspects of the symmetry of n-space:

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– The general linear group GL(n, R) of all real invertible n × n matrices (with

matrix multiplication as operation) contains rotations, reflections, dilations, shear

transformations, and their combinations.

– The orthogonal group O(n, R) of all real orthogonal n × n matrices contains the

rotations and reflections of n-space.

– The special orthogonal group SO(n, R) of all real orthogonal n × n matrices with

determinant 1 contains the rotations of n-space.

**All these matrix groups are Lie groups: groups which are differentiable manifolds such
**

that the group operations are smooth maps.

Other groups

• The trivial group consists only of its identity element.

**• If X is a topological space and x is a point of X, we can define the fundamental group
**

of X at x. It consists of (equivalence classes of) continuous paths starting and ending

at x and describes the structure of the “holes” in X accessible from x.

**• The free groups are important in algebraic topology. In a sense, they are the most
**

general groups, having only those relations among their elements that are absolutely

required by the group axioms.

**• If A and B are two abelian groups (or modules over the same ring), then the set
**

Hom(A, B) of all homomorphisms from A to B is an abelian group (since the sum

and difference of two homomorphisms is again a homomorphism). Note that the

commutativity of B is crucial here: without it, one couldn’t prove that the sum of

two homorphisms is again a homomorphism.

**• The set of all invertible n × n matrices over some ring R forms a group denoted by
**

GL(n, R).

**• The positive integers less than n which are coprime to n form a group if the operation
**

is defined as multiplication modulo n. This is a cyclic group whose order is given by

the Euler phi-function φ(n),

**• Generalizing the last two examples, every ring (and every monoid) contains a group,
**

its group of units (invertible elements), where the group operation is ring (monoid)

multiplication.

**• If K is a number field, then multiplication of (equivalence classes of) non-zero ideals
**

in the ring of algebraic integers OK gives rise to the ideal class group of K.

**• The set of arithmetic functions that take a value other than 0 at 1 form an abelian group
**

under Dirichlet convolution. They include as a subgroup the set of multiplicative functions.

1249

• Consider the curve C = {(x, y) ∈ R2 | y 2 = x3 − x}. Every straight line intersects

this set in three points (counting a point twice if the line is tangent, and allowing for a

point at infinity). If we require that those three points add up to zero for any straight

line, then we have defined an abelian group structure on C. Groups like these are

called abelian varieties; the most prominent examples are elliptic curves of which C is

the simplest one.

• In the classification of all finite simple groups, several “sporadic” groups occur which

don’t follow any discernable pattern. The largest of these is the monster group with

some 8 · 1053 elements.

Version: 14 Owner: AxelBoldt Author(s): AxelBoldt, NeuRet

302.7 group

Group.

A group is a pair (G, ∗) where G is a non-empty set and ∗ is binary operation on G that

holds the following conditions.

**• For any a, b, c ∈ G, (a ∗ b) ∗ c = a ∗ (b ∗ c). (associativity of the operation).
**

• For any a, b in G, a ∗ b belongs to G. (The operation ∗ is closed).

• There is an element e ∈ G such that ge = eg = g for any g ∈ G. (Existence of

identity element).

• For any g ∈ G there exists an element h such that gh = hg = e. (Existence of inverses)

**Usually the symbol ∗ is omitted and we write ab for a ∗ b. Sometimes, the symbol + is used
**

to represent the operation, especially when the group is abelian.

It can be proved that there is only one identity element , and that for every element there

is only one inverse. Because of this we usually denote the inverse of a as a−1 or −a when we

are using additive notation. The identity element is also called neutral element due to its

behavior with respect to the operation.

Version: 10 Owner: drini Author(s): drini

302.8 quotient group

**Let (G, ∗) be a group and H a normal subgroup. The relation ∼ given by a ∼ b when ab−1 ∈
**

H is an equivalence relation. The equivalence classes are called cosets. The equivalence class

of a is denoted as aH (or a + H if additive notation is being used).

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We can induce a group structure on the cosets with the following operation:

(aH) ? (bH) = (a ∗ b)H.

**The collection of cosets is denoted as G/H and together with the ? operation form the
**

quotient group or factor group of G with H.

Example. Consider the group Z and the subgroup

3Z = {n ∈ Z : n = 3k, k ∈ Z}.

Since Z is abelian, 3Z is then also a normal subgroup.

**Using additive notation, the equivalence relation becomes n ∼ m when (n−m) ∈ 3Z, that is,
**

3 divides n − m. So the relation is actually congruence modulo 3. Therefore the equivalence

classes (the cosets) are:

3Z = . . . , −9, −6, −3, 0, 3, 6, 9, . . .

1 + 3Z = . . . , −8, −5, −2, 1, 4, 7, 10, . . .

2 + 3Z = . . . , −7, −4, −1, 2, 5, 8, 11, . . .

which we’ll represent as 0̄, 1̄ and 2̄.

**Then we can check that Z/3Z is actually the integers modulo 3 (that is, Z/3Z ∼
**

= Z3 ).

Version: 6 Owner: drini Author(s): drini

1251

Chapter 303

**20-02 – Research exposition
**

(monographs, survey articles)

303.1 length function

Let G be a group. A length function on G is a function L : G → R+ satisfying:

L(e) = 0,

L(g) = L(g −1 ), ∀g ∈ G,

L(g1 g2 ) 6 L(g1 ) + L(g2 ), ∀g1 , g2 ∈ G.

Version: 2 Owner: mhale Author(s): mhale

1252

Chapter 304

**20-XX – Group theory and
**

generalizations

304.1 free product with amalgamated subgroup

**Definition 26. Let Gk , k = 0, 1, 2 be groups and ik : G0 → Gi , k = 1, 2 be monomorphisms.
**

The free product of G1 and G2 with amalgamated subgroup G0 , is defined to be a group G

that has the following two properties

**1. there are homomorphisms jk : Gk → G, k = 1, 2 that make the following diagram
**

commute

G1

i1 j1

G0 G

i2 j2

G2

**2. G is universal with respect to the previous property, that is for any other group G0 and
**

homomorphisms jk0 : Gk → G0 , k = 1, 2 that fit in such a commutative diagram there

is a unique homomorphism G → G0 so that the following diagram commutes

G1 j10

i1 j1

!

G0 G G0

i2 j2

G2 j20

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It follows by “general nonsense” that the free product of G1 and G2 with amalgamated

subgroup G0 , if it exists, is “unique up to unique isomorphism.” The free product of G1 and

G2 with amalgamated subgroup G0 , is denoted by G1 FG0 G2 . The following theorem asserts

its existence.

**Theorem 1. G1 FG0 G2 exists for any groups Gk , k = 0, 1, 2 and monomorphisms ik : G0 →
**

Gi , k = 1, 2.

**[ Sketch of proof] Without loss of generality assume that G0 is a subgroup of Gk and
**

that ik is the inclusion for k = 1, 2. Let

Gk = h(xk;s )s∈S | (rk;t)t∈T i

**be a presentation of Gk for k = 1, 2. Each g ∈ G0 can be expressed as a word in the generators
**

of Gk ; denote that word by wk (g) and let N be the normal closure of {w1 (g)w2(g)−1 | g ∈ G0 }

in the free product G1 FG2 . Define

G1 FG0 G2 := G1 FG2 /N

**and for k = 0, 1 define jk to be the inclusion into the free product followed by the canonical projection.
**

Clearly (1) is satisfied, while (2) follows from the universal properties of the free product

and the quotient group.

**Notice that in the above proof it would be sufficient to divide by the relations w1 (g)w2 (g)−1
**

for g in a generating set of G0 . This is useful in practice when one is interested in obtaining

a presentation of G1 FG0 G2 .

In case that the ik ’s are not injective the above still goes through verbatim. The group thusly

obtained is called a “pushout”.

Examples of free products with amalgamated subgroups are provided by Van Kampen’s theorem.

Version: 1 Owner: Dr Absentius Author(s): Dr Absentius

304.2 nonabelian group

**Let (G, ∗) be a group. If a ∗ b 6= b ∗ a for some a, b ∈ G, we say that the group is nonabelian
**

or noncommutative.

proposition. There is a nonabelian group for which x 7→ x3 is a homomorphism

Version: 2 Owner: drini Author(s): drini, apmxi

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Chapter 305

**20A05 – Axiomatics and elementary
**

properties

305.1 Feit-Thompson theorem

**An important result in the classification of all finite simple groups, the Feit-Thompson the-
**

orem states that every non-Abelian simple group must have even order.

The proof requires 255 pages.

Version: 1 Owner: mathcam Author(s): mathcam

**305.2 Proof: The orbit of any element of a group is a
**

subgroup

**Following is a proof that, if G is a group and g ∈ G, then hgi ≤ G. Here hgi is the orbit of
**

g and is defined as

hgi = {g n : n ∈ Z}

Since g ∈ hgi, then hgi is nonempty.

**Let a, b ∈ hgi. Then there exist x, y ∈ Z such that a = g x and b = g y . Since ab−1 =
**

g x (g y )−1 = g x g −y = g x−y ∈ hgi, it follows that hgi ≤ G.

Version: 3 Owner: drini Author(s): drini, Wkbj79

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305.3 center

**The center of a group G is the subgroup of elements which commute with every other element.
**

Formally

Z(G) = {x ∈ G | xg = gx, ∀ g ∈ G}

It can be shown that the center has the following properties

• It is non-empty since it contains at least the identity element

• It consists of those conjugacy classes containing just one element

• The center of an abelian group is the entire group

• It is normal in G

• Every p-group has a non-trivial center

Version: 5 Owner: vitriol Author(s): vitriol

305.4 characteristic subgroup

**If (G, ∗) is a group, then H is a characteristic subgroup of G (H char G) if every automorphism
**

of G maps H to itself. That is:

∀f ∈ Aut(G)∀h ∈ Hf (h) ∈ H

or, equivalently:

∀f ∈ Aut(G)f [H] = H

A few properties of characteristic subgroups:

(a) If H char G then H is a normalsubgroup of G

(b) If G has only one subgroup of a given size then that subgroup is characteristic

(c) If K char H and H E G then K E G (contrast with normality of subgroups is not transitive)

(d) If K char H and H char G then K char G

Proofs of these properties:

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(a) Consider H char G under the inner automorphisms of G. Since every automorphism

preserves H, in particular every inner automorphism preserves H, and therefore g ∗

h ∗ g −1 ∈ H for any g ∈ G and h ∈ H. This is precisely the definition of a normal

subgroup.

**(b) Suppose H is the only subgroup of G of order n. In general, homomorphisms takes
**

subgroups to subgroups, and of course isomorphisms take subgroups to subgroups of

the same order. But since there is only one subgroup of G of order n, any automorphism

must take H to H, and so H char G.

**(c) Take K char H and H E G, and consider the inner automorphisms of G (automor-
**

phisms of the form h 7→ g ∗ h ∗ g −1 for some g ∈ G). These all preserve H, and so are

automorphisms of H. But any automorphism of H preserves K, so for any g ∈ G and

k ∈ K, g ∗ k ∗ g −1 ∈ K.

**(d) Let K char H and H char G, and let φ be an automorphism of G. Since H char G,
**

φ[H] = H, so φH , the restriction of φ to H is an automorphism of H. Since K char H,

so φH [K] = K. But φH is just a restriction of φ, so φ[K] = K. Hence K char G.

Version: 1 Owner: Henry Author(s): Henry

305.5 class function

**Given a field K, a K–valued class function on a group G is a function f : G −→ K such
**

that f (g) = f (h) whenever g and h are elements of the same conjugacy class of G.

**An important example of a class function is the character of a group representation. Over
**

the complex numbers, the set of characters of the irreducible representations of G form a

basis for the vector space of all C–valued class functions, when G is a compact Lie group.

**Relation to the convolution algebra Class functions are also known as central func-
**

tions, because they correspond to functions f in the convolution algebra C ∗ (G) that have

the property f ∗ g = g ∗ f for all g ∈ C ∗ (G) (i.e., they commute with everything under the

convolution operation). More precisely, the set of measurable complex valued class func-

tions f is equal to the set of central elements of the convolution algebra C ∗ (G), for G a

locally compact group admitting a Haar measure.

Version: 5 Owner: djao Author(s): djao

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305.6 conjugacy class

**Two elements g and g 0 of a group G are said to be conjugate if there exists h ∈ G such that
**

g 0 = hgh−1 . Conjugacy of elements is an equivalence relation, and the equivalence classes of

G are called conjugacy classes.

Two subsets S and T of G are said to be conjugate if there exists g ∈ G such that

T = {gsg −1 | s ∈ S} ⊂ G.

**In this situation, it is common to write gSg −1 for T to denote the fact that everything in T
**

has the form gsg −1 for some s ∈ S. We say that two subgroups of G are conjugate if they

are conjugate as subsets.

Version: 2 Owner: djao Author(s): djao

305.7 conjugacy class formula

**The conjugacy classes of a group form a partition of its elements. In a finite group, this
**

means that the order of the group is the sum of the number of elements of the distinct

conjugacy classes. For an element g of group G, we denote the conjugacy class of g as Cg

and the normalizer in G of g as NG (g). The number of elements in Cg equals [G : NG (g)], the

index of the normalizer of g in G. For an element g of the center Z(G) of G, the conjugacy

class of g consists of the singleton {g}. Putting this together gives us the conjugacy class

formula m

X

|G| = |Z(G)| + [G : NG (xi )]

i=1

where the xi are elements of the distinct conjugacy classes contained in G − Z(G).

Version: 3 Owner: lieven Author(s): lieven

305.8 conjugate stabilizer subgroups

Let · be a right group action of G on a set M. Then

Gα·g = g −1 Gα g

1

for any α ∈ M and g ∈ G.

Proof:

1

Gα is the stabilizer subgroup of α ∈ M .

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x ∈ Gα·g ↔ α · (gx) = α · g ↔ α · (gxg −1 ) = α ↔ gxg −1 ∈ Gα ↔ x ∈ g −1 αg

and therefore Gα·g = g −1 Gα g.

Thus all stabilizer subgroups for elements of the orbit G(α) of α are conjugate to Gα .

Version: 4 Owner: Thomas Heye Author(s): Thomas Heye

305.9 coset

**Let H be a subgroup of a group G, and let a ∈ G. The left coset of a with respect to H in
**

G is defined to be the set

aH := {ah | h ∈ H}.

The right coset of a with respect to H in G is defined to be the set

Ha := {ha | h ∈ H}.

T

Two left cosets aH and bH of H in G are either identical or disjoint. Indeed, if c ∈ aH bH,

then c = ah1 and c = bh2 for some h1 , h2 ∈ H, whence b−1 a = h2 h−1 1 ∈ H. But then, given

any ah ∈ aH, we have ah = (bb−1 )ah = b(b−1 a)h ∈ bH, so aH ⊂ bH, and similarly

bH ⊂ aH. Therefore aH = bH.

**Similarly, any two right cosets Ha and Hb of H in G are either identical or disjoint. Accord-
**

ingly, the collection of left cosets (or right cosets) partitions the group G; the correspond-

ing equivalence relation for left cosets can be described succintly by the relation a ∼ b if

a−1 b ∈ H, and for right cosets by a ∼ b if ab−1 ∈ H.

**The index of H in G, denoted [G : H], is the cardinality of the set G/H of left cosets of H
**

in G.

Version: 5 Owner: djao Author(s): rmilson, djao

305.10 cyclic group

**A group G is said to be cyclic if it is generated entirely by some x ∈ G. That is, if G has
**

infinite order then every g ∈ G can be expressed as xk with k ∈ Z. If G has finite order

then every g ∈ G can be expressed as xk with k ∈ N0 , and G has exactly φ(|G|) generators,

where φ is the Euler totient function.

**It is a corollary of Lagrange’s theorem that every group of prime order is cyclic. All cyclic
**

groups of the same order are isomorphic to each other. Consequently cyclic groups of order

n are often denoted by Cn . Every cyclic group is abelian.

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Examples of cyclic groups are (Zm , +m ), (Z?p , ×p ) and (Rm , ×m ) where p is prime and Rm =

{n ∈ N : (n, m) = 1, n ≤ m}

Version: 10 Owner: yark Author(s): yark, Larry Hammick, vitriol

305.11 derived subgroup

**Let G be a group and a, b ∈ G. The group element aba−1 b−1 is called the commutator of a
**

and b. An element of G is called a commutator if it is the commutator of some a, b ∈ G.

**The subgroup of G generated by all the commutators in G is called the derived subgroup
**

of G, and also the commutator subgroup. It is commonly denoted by G0 and also by G(1) .

Alternatively, one may define G0 as the smallest subgroup that contains all the commutators.

Note that the commutator of a, b ∈ G is trivial, i.e.

aba−1 b−1 = 1

**if and only if a and b commute. Thus, in a fashion, the derived subgroup measures the degree
**

to which a group fails to be abelian.

Proposition 1. The derived subgroup G0 is normal in G, and the factor group G/G0 is

abelian. Indeed, G is abelian if and only if G0 is the trivial subgroup.

**One can of course form the derived subgroup of the derived subgroup; this is called the
**

second derived subgroup, and denoted by G00 or by G(2) . Proceeding inductively one defines

the nth derived subgroup as the derived subgroup of G(n−1) . In this fashion one obtains a

sequence of subgroups, called the derived series of G:

G = G(0) ⊇ G(1) ⊇ G(2) ⊇ . . .

**Proposition 2. The group G is solvable if and only if the derived series terminates in the
**

trivial group {1} after a finite number of steps. In this case, one can refine the derived series

to obtain a composition series (a.k.a. a Jordan-Holder decomposition) of G.

Version: 4 Owner: rmilson Author(s): rmilson

305.12 equivariant

**Let G be a group, and X and Y left (resp. right) homogeneous spaces of G. Then a map
**

f : X → Y is called equivariant if g(f (x)) = f (gx) (resp. (f (x))g = f (xg)) for all g ∈ G.

Version: 1 Owner: bwebste Author(s): bwebste

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305.13 examples of finite simple groups

**This entry under construction. If I take too long to finish it, nag me about it, or fill in the
**

rest yourself.

All groups considered here are finite.

**It is now widely believed that the classification of all finite simple groups up to isomorphism
**

is finished. The proof runs for at least 10,000 printed pages, and as of the writing of this

entry, has not yet been published in its entirety.

Abelian groups

**• The first trivial example of simple groups are the cyclic groups of prime order. It is
**

not difficult to see (say, by Cauchy’s theorem) that these are the only abelian simple

groups.

Alternating groups

**• The alternating group on n symbols is the set of all even permutations of Sn , the
**

symmetric group on n symbols. It is usually denoted by An , or sometimes by Alt(n).

This is a normal subgroup of Sn , namely the kernel of the homomorphism that sends

every even permutation to 1 and the odd permutations to −1. Because every per-

mutation is either even or odd, and there is a bijection between the two (multiply

every even permutation by a transposition), the index of An in Sn is 2. A3 is simple

because it only has three elements, and the simplicity of An for n > 5 can be proved

by an elementary argument. The simplicity of the alternating groups is an important

fact that Évariste Galois required in order to prove the insolubility by radicals of the

general polynomial of degree higher than four.

Groups of Lie type

• Projective special linear groups

• Other groups of Lie type.

**Sporadic groups There are twenty-six sporadic groups (no more, no less!) that do not
**

fit into any of the infinite sequences of simple groups considered above. These often arise as

the group of automorphisms of strongly regular graphs.

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• Mathieu groups.

• Janko groups.

• The baby monster.

• The monster.

Version: 8 Owner: drini Author(s): bbukh, yark, NeuRet

305.14 finitely generated group

**A group G is finitely generated if there is a finite subset X ⊆ G such that X generates G.
**

That is, every element of G is a product of elements of X and inverses of elements of X. Or,

equivalently, no proper subgroup of G contains X.

**Every finite group is finitely generated, as we can take X = G. Every finitely generated
**

group is countable.

Version: 6 Owner: yark Author(s): yark, nerdy2

305.15 first isomorphism theorem

**If f : G → H is a homorphism of groups (or rings, or modules), then it induces an
**

isomorphism G/ker f ≈ imf .

Version: 2 Owner: nerdy2 Author(s): nerdy2

305.16 fourth isomorphism theorem

fourth isomorphism theorem

1: X group

2: N E X

3: A set of subgroups of X that contain N

4: B set of subgroups of X/N

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5: ∃ϕ : A → B bijection : ∀Y, Z 6 X : *N 6 Y & N 6 Z+ ⇒ **Y 6 Z ⇔ Y /N T 6 Z/N +

& *Z T6 Y ⇒ [Y : Z] = [Y /N : Z/N] + & *hY, Zi/N = hY /N, Z/Ni + & *(Y Z)/N =

Y /N Z/N + & *Y E G ⇔ Y /N E G/N+

Note: This is a “seed” entry written using a short-hand format described in this FAQ.

Version: 2 Owner: bwebste Author(s): yark, apmxi

305.17 generator

If G is a cyclic group and g ∈ G, then g is a generator of G if hgi = G.

**All infinite cyclic groups have exactly 2 generators. Let G be an infinite cyclic group and
**

g be a generator of G. Let z ∈ Z such that g z is a generator of G. Then hg z i = G. Since

g ∈ G, then g ∈ hg z i. Thus, there exists, n ∈ Z with g = (g z )n = g nz . Thus, g nz−1 = eG .

Since G is infinite and |g| = |hgi| = |G| must be infinity, then nz − 1 = 0. Since nz = 1 and

n and z are integers, then n = z = 1 or n = z = −1. It follows that the only generators of

G are g and g −1.

**A finite cyclic group of order n has exactly ϕ(n) generators, where ϕ is the Euler totient function.
**

Let G be a finite cyclic group of order n and g be a generator of G. Then |g| = |hgi| = |G| = n.

Let z ∈ Z such that g z is a generator of G. By the division algorithm, there exist q, r ∈ Z with

0 ≤ r < n such that z = qn + r. Thus, g z = g qn+r = g qn g r = (g n )q g r = (eG )q g r = eG g r = g r .

|g| n

Since g r is a generator of G, then hg r i = G. Thus, n = |G| = |hg r i| = |g r | = gcd(r,|g|) = gcd(r,n) .

Therefore, gcd(r, n) = 1, and the result follows.

Version: 3 Owner: Wkbj79 Author(s): Wkbj79

305.18 group actions and homomorphisms

**Notes on group actions and homomorphisms Let G be a group, X a non-empty set and SX
**

the symmetric group of X, i.e. the group of all bijective maps on X. · may denote a left

group action of G on X.

**1. For each g ∈ G we define
**

fg : X −→ X, fg (x) = g · x ∀ x ∈ X.

fg− 1 (fg (x)) = g −1 · (g · x) = x ∀ x ∈ X, so fg− 1 is the inverse of fg . so fg is bijective and

thus element of SX . We define F : G −→ SX , F (g) = fg for all g ∈ G. This mapping

is a group homomorphism: Let g, h ∈ G, x ∈ X. Then

F (gh)(x) = fgh (x) = (gh) · x = g · (h · x) = (fg ◦ fh )(x) =

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implies F (gh) = F (g) ◦ F (h). – The same is obviously true for a right group action.

2. Now let F : G −→ Sx be a group homomorphism, and let f : G × X −→ X, (g, x) −→

F (g)(x) satisfies

(a) f (1G , x) = F (1g )(x) = x∀x ∈ X and

(b) f (gh, x) = F (gh)(x) = (F (g) ◦ F (h)(x) = F (g)(F (h)(x)) = f (g, f (h, x)),

so f is a group action induced by F .

Characterization of group actions

**Let G be a group acting on a set X. Using the same notation as above, we have for each
**

g ∈ ker(F ) [

F (g) = id x = fg ↔ g · x = x ∀x ∈ X ↔ g ∈ Gx (305.18.1)

x∈X

and it follows [

ker(F ) = Gx .

x∈X

**Let G act transitively on X. Then for any x ∈ X X is the orbit G(x) of x. As shown in
**

“conjugate stabilizer subgroups’, all stabilizer subgroups of elements y ∈ G(x) are conjugate subgroups

to Gx in G. From the above it follows that

[

ker(F ) = gGx g −1 .

g∈G

**For a faithful operation of G the condition g · x = x∀x ∈ X → g = 1G is equivalent to
**

ker(F ) = {1G }

and therefore F : G −→ SX is a monomorphism.

**For the trivial operation of G on X given by g · x = x∀g ∈ G the stabilizer subgroup Gx is
**

G for all x ∈ X, and thus

ker(F ) = G.

The corresponding homomorphism is g −→ id x∀g ∈ G.

**If the operation of G on X is free, then Gx = {1G } ∀ x ∈ X, thus the kernel of F is {1G }–like
**

for a faithful operation. But:

**Let X = {1, . . . , n} and G = Sn . Then the operation of G on X given by
**

π · i := π(i)∀i ∈ X, π ∈ Sn

is faithful but not free.

Version: 5 Owner: Thomas Heye Author(s): Thomas Heye

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305.19 group homomorphism

**Let (G, ∗g ) and (K, ∗k ) be two groups. A group homomorphism is a function φ : G → K
**

such that φ(s ∗g t) = φ(s) ∗k φ(t) for all s, t ∈ G.

The composition of group homomorphisms is again a homomorphism.

Let φ : G → K a group homomorphism. Then

**• φ(eg ) = ek where eg and ek are the respective identity elements for G and K.
**

• φ(g)−1 = φ(g −1 ) for all g ∈ G

• φ(g)z = φ(g z ) for all g ∈ G and for all z ∈ Z

The kernel of φ is a subgroup of G and its image is a subgroup of K.

**Some special homomorphisms have special names. If φ : G → K is injective, we say that φ is
**

an monomorphism, and if φ is onto we call it an epimorphism. When φ is both injective and

surjective (that is, bijective) we call it an isomorphism. In the latter case we also say that G

and K are isomorphic, meaning they are basically the same group (have the same structure).

An homomorphism from G on itself is called an endomorphism, and if it is bijective, then is

called an automorphism.

Version: 15 Owner: drini Author(s): saforres, drini

305.20 homogeneous space

**Overview and definition. Let G be a group acting transitively on a set X. In other
**

words, we consider a homomorphism φ : G → Perm(X), where the latter denotes the group

of all bijections of X. If we consider G as being, in some sense, the automorphisms of X,

the transitivity assumption means that it is impossible to distinguish a particular element

of X from any another element. Since the elements of X are indistinguishable, we call X a

homogeneous space. Indeed, the concept of a homogeneous space, is logically equivalent

to the concept of a transitive group action.

**Action on cosets. Let G be a group, H < G a subgroup, and let G/H denote the set of
**

left cosets, as above. For every g ∈ G we consider the mapping ψH (g) : G/H → G/H with

action

aH → gaH, a ∈ G.

Proposition 3. The mapping ψH (g) is a bijection. The corresponding mapping ψH : G →

Perm(G/H) is a group homomorphism, specifying a transitive group action of G on G/H.

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Thus, G/H has the natural structure of a homogeneous space. Indeed, we shall see that

every homogeneous space X is isomorphic to G/H, for some subgroup H.

**N.B. In geometric applications, the want the homogeneous space X to have some extra
**

structure, like a topology or a differential structure. Correspondingly, the group of automor-

phisms is either a continuous group or a Lie group. In order for the quotient space X to

have a Hausdorff topology, we need to assume that the subgroup H is closed in G.

**The isotropy subgroup and the basepoint identification. Let X be a homogeneous
**

space. For x ∈ X, the subgroup

Hx = {h ∈ G : hx = x},

consisting of all G-actions that fix x, is called the isotropy subgroup at the basepoint x. We

identify the space of cosets G/Hx with the homogeneous space by means of the mapping

τx : G/Hx → X, defined by

τx (aHx ) = ax, a ∈ G.

Proposition 4. The above mapping is a well-defined bijection.

**To show that τx is well defined, let a, b ∈ G be members of the same left coset, i.e. there
**

exists an h ∈ Hx such that b = ah. Consequently

bx = a(hx) = ax,

as desired. The mapping τx is onto because the action of G on X is assumed to be transitive.

To show that τx is one-to-one, consider two cosets aHx , bHx , a, b ∈ G such that ax = bx.

It follows that b−1 a fixes x, and hence is an element of Hx . Therefore aHx and bHx are the

same coset.

**The homogeneous space as a quotient. Next, let us show that τx is equivariant relative
**

to the action of G on X and the action of G on the quotient G/Hx .

Proposition 5. We have that

φ(g) ◦ τx = τx ◦ ψHx (g)

for all g ∈ G.

**To prove this, let g, a ∈ G be given, and note that
**

ψHx (g)(aHx ) = gaHx .

The latter coset corresponds under τx to the point gax, as desired.

**Finally, let us note that τx identifies the point x ∈ X with the coset of the identity element
**

eHx , that is to say, with the subgroup Hx itself. For this reason, the point x is often called

the basepoint of the identification τx : G/Hx → X.

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The choice of basepoint. Next, we consider the effect of the choice of basepoint on the

quotient structure of a homogeneous space. Let X be a homogeneous space.

Proposition 6. The set of all isotropy subgroups {Hx : x ∈ X} forms a single conjugacy class

of subgroups in G.

**To show this, let x0 , x1 ∈ X be given. By the transitivity of the action we may choose a
**

ĝ ∈ G such that x1 = ĝx0 . Hence, for all h ∈ G satisfying hx0 = x0 , we have

(ĝhĝ −1 )x1 = ĝ(h(ĝ −1 x1 )) = ĝx0 = x1 .

Similarly, for all h ∈ Hx1 we have that ĝ −1hĝ fixes x0 . Therefore,

ĝ(Hx0 )ĝ −1 = Hx1 ;

or what is equivalent, for all x ∈ X and g ∈ G we have

gHx g −1 = Hgx .

**Equivariance. Since we can identify a homogeneous space X with G/Hx for every possible
**

x ∈ X, it stands to reason that there exist equivariant bijections between the different G/Hx .

To describe these, let H0 , H1 < G be conjugate subgroups with

H1 = ĝH0 ĝ −1

for some fixed ĝ ∈ G. Let us set

X = G/H0 ,

and let x0 denote the identity coset H0 , and x1 the coset ĝH0 . What is the subgroup of G

that fixes x1 ? In other words, what are all the h ∈ G such that

hĝH0 = ĝH0 ,

or what is equivalent, all h ∈ G such that

ĝ −1 hĝ ∈ H0 .

The collection of all such h is precisely the subgroup H1 . Hence, τx1 : G/H1 → G/H0 is the

desired equivariant bijection. This is a well defined mapping from the set of H1 -cosets to

the set of H0 -cosets, with action given by

τx1 (aH1 ) = aĝH0 , a ∈ G.

**Let ψ0 : G → Perm(G/H0 ) and ψ1 : G → Perm(G/H1 ) denote the corresponding coset
**

G-actions.

Proposition 7. For all g ∈ G we have that

τx1 ◦ ψ1 (g) = ψ0 (g) ◦ τx1 .

Version: 3 Owner: rmilson Author(s): rmilson

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305.21 identity element

**Let G be a groupoid, that is a set with a binary operation G×G → G, written muliplicatively
**

so that (x, y) 7→ xy.

An identity element for G is an element e such that ge = eg = g for all g ∈ G.

**The symbol e is most commonly used for identity elements. Another common symbol for
**

an identity element is 1, particularly in semigroup theory (and ring theory, considering the

multiplicative structure as a semigroup).

**Groups, monoids, and loops are classes of groupoids that, by definition, always have an
**

identity element.

Version: 6 Owner: mclase Author(s): mclase, vypertd, imran

305.22 inner automorphism

Let G be a group. For every x ∈ G, we define a mapping

φx : G → G, y 7→ xyx−1 , y ∈ G,

called conjugation by x. It is easy to show the conjugation map is in fact, a group automorphism.

**An automorphism of G that corresponds to the conjugation by some x ∈ G is called inner.
**

An automorphism that isn’t inner is called an outer automorphism.

**The composition operation gives the set of all automorphisms of G the structure of a group,
**

Aut(G). The inner automorphisms also form a group, Inn(G), which is a normal subgroup

of Aut(G). Indeed, if φx , x ∈ G is an inner automorphism and π : G → G an arbitrary

automorphism, then

π ◦ φx ◦ π −1 = φπ(x) .

Let us also note that the mapping

x 7→ φx , x∈G

**is a surjective group homomorphism with kernel Z(G), the centre subgroup. Consequently,
**

Inn(G) is naturally isomorphic to the quotient of G/ Z(G).

Version: 7 Owner: rmilson Author(s): rmilson, tensorking

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305.23 kernel

**Let ρ : G → K be a group homomorphism. The preimage of the codomain identity element
**

eK ∈ K forms a subgroup of the domain G, called the kernel of the homomorphism;

ker(ρ) = {s ∈ G | ρ(s) = eK }

The kernel is a normal subgroup. It is the trivial subgroup if and only if ρ is a monomorphism.

Version: 9 Owner: rmilson Author(s): rmilson, Daume

305.24 maximal

**Let G be a group. A subgroup H of G is said to be maximal if H 6= G and whenever K is
**

a subgroup of G with H ⊆ K ⊆ G then K = H or K = G.

Version: 1 Owner: Evandar Author(s): Evandar

305.25 normal subgroup

**A subgroup H of a group G is normal if aH = Ha for all a ∈ G. Equivalently, H ⊂ G is
**

normal if and only if aHa−1 = H for all a ∈ G, i.e., if and only if each conjugacy class of G

is either entirely inside H or entirely outside H.

The notation H E G or H / G is often used to denote that H is a normal subgroup of G.

**The kernel ker (f ) of any group homomorphism f : G −→ G0 is a normal subgroup of G.
**

More surprisingly, the converse is also true: any normal subgroup H ⊂ G is the kernel of

some homomorphism (one of these being the projection map ρ : G −→ G/H, where G/H is

the quotient group).

Version: 6 Owner: djao Author(s): djao

305.26 normality of subgroups is not transitive

**Let G be a group. Obviously, a subgroup K ≤ H of a subgroup H ≤ G of G is a subgroup
**

K ≤ G of G. It seems plausible that a similar situation would also hold for normal subgroups.

**This is not true. Even when K E H and H E G, it is possible that K 5 G. Here are two
**

examples:

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1. Let G be the subgroup of orientation-preserving isometries of the plane R2 (G is just

all rotations and translations), let H be the subgroup of G of translations, and let K

be the subgroup of H of integer translations τi,j (x, y) = (x + i, y + j), where i, j ∈ Z.

Any element g ∈ G may be represented as g = r1 ◦ t1 = t2 ◦ r2 , where r1,2 are rotations

and t1,2 are translations. So for any translation t ∈ H we may write

g −1 ◦ t ◦ g = r −1 ◦ t0 ◦ r,

**where t0 ∈ H is some other translation and r is some rotation. But this is an orientation-
**

preserving isometry of the plane that does not rotate, so it too must be a translation.

Thus G−1 HG = H, and H E G.

H is an abelian group, so all its subgroups, K included, are normal.

We claim that K 5 G. Indeed, if ρ ∈ G is rotation by 45◦ about the origin, then

ρ−1 ◦ τ1,0 ◦ ρ is not an integer translation.

**2. A related example uses finite subgroups. Let G = D4 be the dihedral group with four
**

elements (the group of automorphisms of the graph of the square C4 ). Then

D4 = r, f | f 2 = 1, r 4 = 1, f r = r −1 f

is generated by r, rotation, and f , flipping.

The subgroup

H = hrf, f ri = 1, rf, r 2 , f r ∼= C2 × C2

is isomorphic to the Klein 4-group – an identity and 3 elements of order 2. H E G

since [G : H] = 2. Finally, take

K = hrf i = {1, rf } E H.

We claim that K 5 G. And indeed,

f ◦ rf ◦ f = f r ∈

/ K.

Version: 4 Owner: ariels Author(s): ariels

305.27 normalizer

Let G be a group, and let H ⊆ G. The normalizer of H in G, written NG (H), is the set

{g ∈ G | gHg −1 = H}

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If H is a subgroup of G, then NG (H) is a subgroup of G containing H.

**Note that H is a normal subgroup of NG (H); in fact, NG (H) is the largest subgroup of
**

G of which H is a normal subgroup. In particular, if H is a normal subgroup of G, then

NG (H) = G.

Version: 6 Owner: saforres Author(s): saforres

305.28 order (of a group)

**The order of a group G is the number of elements of G, denoted |G|; if |G| is finite, then G
**

is said to be a finite group.

**The order of an element g ∈ G is the smallest positive integer n such that g n = e, where e
**

is the identity element; if there is no such n, then g is said to be of infinite order.

Version: 5 Owner: saforres Author(s): saforres

305.29 presentation of a group

**A presentation of a group G is a description of G in terms of generators and relations. We
**

say that the group is finitely presented, if it can be described in terms of a finite number

of generators and a finite number of defining relations. A collection of group elements

gi ∈ G, i ∈ I is said to generate G if every element of G can be specified as a product of the

gi , and of their inverses. A relation is a word over the alphabet consisting of the generators

gi and their inverses, with the property that it multiplies out to the identity in G. A set of

relations rj , j ∈ J is said to be defining, if all relations in G can be given as a product of

the rj , their inverses, and the G-conjugates of these.

**The standard notation for the presentation of a group is
**

G = hgi | rj i,

meaning that G is generated by generators gi , subject to relations rj . Equivalently, one has

a short exact sequence of groups

1 → N → F [I] → G → 1,

where F [I] denotes the free group generated by the gi , and where N is the smallest normal subgroup

containing all the rj . By the Nielsen-Schreier theorem, the kernel N is itself a free group, and

hence we assume without loss of generality that there are no relations among the relations.

**Example. The symmetric group on n elements 1, . . . , n admits the following finite pre-
**

sentation (Note: this presentation is not canonical. Other presentations are known.) As

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generators take

gi = (i, i + 1), i = 1, . . . , n − 1,

the transpositions of adjacent elements. As defining relations take

(gi gj )ni,j = id, i, j = 1, . . . n,

where

ni,i = 1

ni,i+1 = 3

ni,j = 2, i<j+1

This means that a finite symmetric group is a Coxeter group,

Version: 11 Owner: rmilson Author(s): rmilson

305.30 proof of first isomorphism theorem

**Let K denote ker f . K is a normal subgroup of G because, by the following calculation,
**

gkg −1 ∈ K for all g ∈ G and k ∈ K (rules of homomorphism imply the first equality,

definition of K for the second):

f (gkg −1) = f (g)f (k)f (g)−1 = f (g)1H f (g)−1 = 1H

Therefore, G/K is well defined.

**Define a group homomorphism θ : G/K → imf given by:
**

θ(gK) = f (g)

We argue that θ is an isomorphism.

**First, θ is well defined. Take two representative, g1 and g2 , of the same modulo class. By
**

definition, g1 g2−1 is in K. Hence, f sends g1 g2−1 to 1 (all elements of K are sent by f to

1). Consequently, the next calculation is valid: f (g1)f (g2 )−1 = f (g1 g2−1) = 1 but this is the

same as saying that f (g1 ) = f (g2 ). And we are done because the last equality indicate that

θ(g1 K) is equal to θ(g2 K).

**Going backward the last argument, we get that θ is also an injection: If θ(g1 K) is equal to
**

θ(g2 K) then f (g1 ) = f (g2 ) and hence g1 g2−1 ∈ K (exactly as in previous part) which implies

an equality between g1 K and g2 K.

**Now, θ is a homomorphism. We need to show that θ(g1 K · g2 K) = θ(g1 K)θ(g2 K) and that
**

θ((gK)−1 ) = (θ(gK))−1 . And indeed:

θ(g1 K · g2 K) = θ(g1 g2 K) = f (g1 g2 ) = f (g1 )f (g2 ) = θ(g1 K)θ(g2 K)

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θ((gK)−1 ) = θ(g −1 K) = f (g −1 ) = (f (g))−1 = (θ(gK))−1

**To conclude, θ is surjective. Take h to be an element of imf and g its pre-image. Since
**

h = f (g) we have that h is also the image of of θ(gK).

Version: 3 Owner: uriw Author(s): uriw

305.31 proof of second isomorphism theorem

**First, we shall prove that HK is a subgroup of G: Since e ∈ H and e ∈ K, clearly e = e2 ∈
**

HK. Take h1 , h2 ∈ H, k1, k2 ∈ K. Clearly h1 k1 , h2 k2 ∈ HK. Further,

h1 k1 h2 k2 = h1 (h2 h−1 −1

2 )k1 h2 k2 = h1 h2 (h2 k1 h2 )k2

**Since K is a normal subgroup of G and h2 ∈ G, then h−1 −1
**

2 k1 h2 ∈ K. Therefore h1 h2 (h2 k1 h2 )k2 ∈

HK, so HK is closed under multiplication.

Also, (hk)−1 ∈ HK for h ∈ H, k ∈ K, since

(hk)−1 = k −1 h−1 = h−1 hk −1 h−1

**and hk −1 h−1 ∈ K since K is a normal subgroup of G. So HK is closed under inverses, and
**

is thus a subgroup of G.

**Since HK is a subgroup of G, the normality of K in HK follows immediately from the
**

normality of K in G.

T

Clearly H K is a subgroup of G, since it is the intersection of two subgroups of G.

**Finally, define φ : H → HK/K by ϕ(h) = hK. We claim that φ is a surjective homomorphism
**

from H to HK/K. Let h0 k0 K be some element of HK/K; since k0 ∈ K, then h0 k0 K = h0 K,

and φ(h0 ) = h0 K. Now

ker (φ) = {h ∈ H | φ(h) = K} = {h ∈ H | hK = K}

**and if hK = K, then we must have h ∈ K. So
**

\

ker (φ) = {h ∈ H | h ∈ K} = H K

T

Thus, since

T φ(H) = HK/K and ker φ = H K, by the first isomorphism theorem we

T see

that H K is normal in H and that there is a natural isomorphism between H/(H K)

and HK/K.

Version: 8 Owner: saforres Author(s): saforres

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305.32 proof that all cyclic groups are abelian

Following is a proof that all cyclic groups are abelian.

**Let G be a cyclic group and g be a generator of G. Let a, b ∈ G. Then there exist x, y ∈ Z
**

such that a = g x and b = g y . Since ab = g x g y = g x+y = g y+x = g y g x = ba, it follows that G

is abelian.

Version: 2 Owner: Wkbj79 Author(s): Wkbj79

**305.33 proof that all cyclic groups of the same order
**

are isomorphic to each other

The following is a proof that all cyclic groups of the same order are isomorphic to each other.

**Let G be a cyclic group and g be a generator of G. Define ϕ : Z → G by ϕ(c) = g c . Since
**

ϕ(a + b) = g a+b = g a g b = ϕ(a)ϕ(b), then ϕ is a group homomorphism. If h ∈ G, then there

exists x ∈ Z such that h = g x . Since ϕ(x) = g x = h, then ϕ is surjective.

ker ϕ = {c ∈ Z|ϕ(c) = eG } = {c ∈ Z|g c = eG }

**If G is infinite, then ker ϕ = {0}, and ϕ is injective. Hence, ϕ is a group isomorphism, and
**

G∼ = Z.

**If G is finite, then let |G| = n. Thus, |g| = |hgi| = |G| = n. If g c = eG , then n divides c.
**

Therefore, ker ϕ = nZ. By the first isomorphism theorem, G ∼ Z ∼

= nZ = Zn .

**Let H and K be cyclic groups of the same order. If H and K are infinite, then, by the
**

above argument, H ∼

= Z and K ∼ = Z. If H and K are finite of order n, then, by the above

argument, H = Zn and K ∼

∼ = Zn . In any case, it follows that H ∼

= K.

Version: 1 Owner: Wkbj79 Author(s): Wkbj79

**305.34 proof that all subgroups of a cyclic group are
**

cyclic

Following is a proof that all subgroups of a cyclic group are cyclic.

**Let G be a cyclic group and H ≤ G. If G is trivial, then H = G, and H is cyclic. If H is
**

the trivial subgroup, then H = {eG } = heG i, and H is cyclic. Thus, for the remainder of the

proof, it will be assumed that both G and H are nontrivial.

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Let g be a generator of G. Let n be the smallest positive integer such that g n ∈ H.

Claim: H = hg n i

**Let a ∈ hg n i. Then there exists z ∈ Z with a = (g n )z . Since g n ∈ H, then (g n )z ∈ H. Thus,
**

a ∈ H. Hence, hg n i ⊆ H.

**Let h ∈ H. Then h ∈ G. Let x ∈ Z with h = g x . By the division algorithm, there exist
**

q, r ∈ Z with 0 ≤ r < n such that x = qn + r. Since h = g x = g qn+r = g qn g r = (g n )q g r , then

g r = h(g n )−q . Since h, g n ∈ H, then g r ∈ H. By choice of n, r cannot be positive. Thus,

r = 0. Therefore, h = (g n )q g 0 = (g n )q eG = (g n )q ∈ hg n i. Hence, H ⊆ hg n i.

Since hg n i ⊆ H and H ⊆ hg n i, then H = hg n i. It follows that every subgroup of G is cyclic.

Version: 3 Owner: Wkbj79 Author(s): Wkbj79

305.35 regular group action

**Let G be a group action on a set X. The action is called regular if for any pair α, β ∈ X
**

there exists exactly one g ∈ G such that g · α = β. (For a right group action it is defined

correspondingly.)

Version: 3 Owner: Thomas Heye Author(s): Thomas Heye

305.36 second isomorphism theorem

**Let (G, ∗) be a group. Let H be a subgroup of G and let K be a normal subgroup of G.
**

Then

• HK := {h ∗ k | h ∈ H, k ∈ K} is a subgroup of G,

• K is a normal subgroup of HK,

T

• H K is a normal subgroup of H,

T

• There is a natural group isomorphism H/(H K) = HK/K.

**The same statement also holds in the category of modules over a fixed ring (where normality is
**

neither needed nor relevant), and indeed can be formulated so as to hold in any abelian category.

Version: 4 Owner: djao Author(s): djao

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305.37 simple group

**Let G be a group. G is said to be simple if the only normal subgroups of G are {1} and G
**

itself.

Version: 3 Owner: Evandar Author(s): Evandar

305.38 solvable group

A group G is solvable if it has a composition series

G = G0 ⊃ G1 ⊃ · · · ⊃ Gn = {1}

where all the quotient groups Gi /Gi+1 are abelian.

Version: 4 Owner: djao Author(s): djao

305.39 subgroup

Definition:

Let (G, ∗) be a group and let K be subset of G. Then K is a subgroup of G defined under

the same operation if K is a group by itself (respect to ∗), that is:

• K is closed under the ∗ operation.

• There exists an identity element e ∈ K such that for all k ∈ K, k ∗ e = k = e ∗ k.

• Let k ∈ K then there exists an inverse k −1 ∈ K such that k −1 ∗ k = e = k ∗ k −1 .

**The subgroup is denoted likewise (K, ∗). We denote K being a subgroup of G by writing
**

K 6 G.

properties:

**• The set {e} whose only element is the identity is a subgroup of any group. It is called
**

the trivial subgroup.

• Every group is a subgroup of itself.

**• The null set {} is never a subgroup (since the definition of group states that the set
**

must be non-empty).

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There is a very useful theorem that allows proving a given subset is a subgroup.

Theorem:

If K is a nonempty subset of of the group G. Then K is a subgroup of G if and only if

s, t ∈ K implies that st−1 ∈ K.

**Proof: First we need to show if K is a subgroup of G then st−1 ∈ K. Since s, t ∈ K then
**

st−1 ∈ K, because K is a group by itself.

Now, suppose that if for any s, t ∈ K ⊆ G we have st−1 ∈ K. We want to show that K is a

subgroup, which we will acomplish by proving it holds the group axioms.

**Since tt−1 ∈ K by hypothesis, we conclude that the identity element is in K: e ∈ K.
**

(Existence of identity)

**Now that we know e ∈ K, for all t in K we have that et−1 = t−1 ∈ K so the inverses of
**

elements in K are also in K. (Existence of inverses).

**Let s, t ∈ K. Then we know that t−1 ∈ K by last step. Applying hypothesis shows that
**

s(t−1 )−1 = st ∈ K

so K is closed under the operation. QED

Example:

**• Consider the group (Z, +). Show that(2Z, +) is a subgroup.
**

The subgroup is closed under addition since the sum of even integers is even.

The identity 0 of Z is also on 2Z since 2 divides 0. For every k ∈ 2Z there is an

−k ∈ 2Z which is the inverse under addition and satisfies −k + k = 0 = k(−k).

Therefore (2Z, +) is a subgroup of (Z, +).

Another way to show (2Z, +) is a subgroup is by using the proposition stated above. If

s, t ∈ 2Z then s, t are even numbers and s −t ∈ 2Z since the difference of even numbers

is always an even number.

See also:

• Wikipedia, subgroup

Version: 7 Owner: Daume Author(s): Daume

305.40 third isomorphism theorem

**If G is a group (or ring, or module) and H ⊂ K are normal subgroups (or ideals, or submod-
**

ules), with H normal (or an ideal, or a submodule) in K then there is a natural isomorphism

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(G/H)/(K/H) ≈ G/K.

I think it is not uncommon to see the third and second isomorphism theorems permuted.

Version: 2 Owner: nerdy2 Author(s): nerdy2

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Chapter 306

20A99 – Miscellaneous

306.1 Cayley table

**A Cayley table for a group is essentially the “multiplication table” of the group.1 The
**

columns and rows of the table (or matrix) are labeled with the elements of the group, and

the cells represent the result of applying the group operation to the row-th and column-th

elements.

**Formally, Let G be our group, with operation ◦ the group operation. Let C be the Cayley
**

table for the group, with C(i, j) denoting the element at row i and column j. Then

C(i, j) = ei ◦ ej

where ei is the ith element of the group, and ej is the jth element.

**Note that for an abelian group, we have ei ◦ ej = ej ◦ ei , hence the Cayley table is a
**

symmetric matrix.

**All Cayley tables for isomorphic groups are isomorphic (that is, the same, invariant of the
**

labeling and ordering of group elements).

306.1.1 Examples.

• The Cayley table for Z4 , the group of integers modulo 4 (under addition), would be

1

A caveat to novices in group theory: multiplication is usually used notationally to represent the group

operation, but the operation needn’t resemble multiplication in the reals. Hence, you should take “multipli-

cation table” with a grain or two of salt.

1279

[0] [1] [2] [3]

[0] [0] [1] [2] [3]

[1] [1] [2] [3] [0]

[2] [2] [3] [0] [1]

[3] [3] [0] [1] [2]

**• The Cayley table for S3 , the permutation group of order 3, is
**

(1) (123) (132) (12) (13) (23)

(1) (1) (123) (132) (12) (13) (23)

23) (123) (132) (1) (13) (23) (12)

(132) (132) (1) (123) (23) (12) (13)

2) (12) (23) (13) (1) (132) (123)

(13) (13) (12) (23) (123) (1) (132)

(23) (23) (13) (12) (132) (123) (1)

Version: 6 Owner: akrowne Author(s): akrowne

306.2 proper subgroup

A group H is a proper subgroup of a group G if and only if H is a subgroup of G and

H 6= G. (306.2.1)

**Similarly a group H is an improper subgroup of a group G if and only if H is a subgroup
**

of G and

H = G. (306.2.2)

Version: 2 Owner: imran Author(s): imran

306.3 quaternion group

**The quaternion group, or quaternionic group, is a noncommutative group with eight ele-
**

ments. It is traditionally denoted by Q (not to be confused with Q) or by Q8 . This group

is defined by the presentation

{i, j; i4 , i2 j 2 , iji−1 j −1 }

or, equivalently, defined by the multiplication table

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· 1 i j k −i −j −k −1

1 1 i j k −i −j −k −1

i i −1 k −j 1 −k j −i

j j −k −1 i k 1 −i −j

k k j −i −1 −j i 1 −k

−i −i 1 −k j −1 k −j i

−j −j k 1 −i −k −1 i j

−k −k −j i 1 j −i −1 k

−1 −1 −i −j −k i j k 1

where we have put each product xy into row x and column y. The minus signs are justified

by the fact that {1, −1} is subgroup contained in the center of Q. Every subgroup of Q is

normal and, except for the trivial subgroup {1}, contains {1, −1}. The dihedral group D4

(the group of symmetries of a square) is the only other noncommutative group of order 8.

**Since i2 = j 2 = k 2 = −1, the elements i, j, and k are known as the imaginary units, by
**

analogy with i ∈ C. Any pair of the imaginary units generate the group. Better, given

x, y ∈ {i, j, k}, any element of Q is expressible in the form xm y n .

**Q is identified with the group of units (invertible elements) of the ring of quaternions over
**

Z. That ring is not identical to the group ring Z[Q], which has dimension 8 (not 4) over Z.

Likewise the usual quaternion algebra is not quite the same thing as the group algebra R[Q].

**Quaternions were known to Gauss in 1819 or 1820, but he did not publicize this discovery,
**

and quaternions weren’t rediscovered until 1843, with Hamilton. For an excellent account of

this famous Story, see http://math.ucr.edu/home/baez/Octonions/node1.html.

**Version: 6 Owner: vernondalhart Author(s): vernondalhart, Larry Hammick, patrickwon-
**

ders

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Chapter 307

**20B05 – General theory for finite
**

groups

307.1 cycle notation

**The cycle notation is a useful convention for writing down a permutations in terms of its
**

constituent cycles. Let S be a finite set, and

a1 , . . . , ak , k>2

distinct elements of S. The expression (a1 , . . . , ak ) denotes the cycle whose action is

a1 7→ a2 7→ a3 . . . ak 7→ a1 .

Note there are k different expressions for the same cycle; the following all represent the same

cycle:

(a1 , a2 , a3 , . . . , ak ) = (a2 , a3 , . . . , ak , a1 ), = . . . = (ak , a1 , a2 , . . . , ak−1 ).

Also note that a 1-element cycle is the same thing as the identity permutation, and thus

there is not much point in writing down such things. Rather, it is customary to express the

identity permutation simply as ().

**Let π be a permutation of S, and let
**

S1 , . . . , Sk ⊂ S, k∈N

be the orbits of π with more than 1 element. For each j = 1, . . . , k let nj denote the

cardinality of Sj . Also, choose an a1,j ∈ Sj , and define

ai+1,j = π(ai,j ), i ∈ N.

We can now express π as a product of disjoint cycles, namely

π = (a1,1 , . . . an1 ,1 )(a2,1 , . . . , an2 ,2 ) . . . (ak,1 , . . . , ank ,k ).

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By way of illustration, here are the 24 elements of the symmetric group on {1, 2, 3, 4} ex-

pressed using the cycle notation, and grouped according to their conjugacy classes:

(),

(12), (13), (14), (23), (24), (34)

(123), (213), (124), (214), (134), (143), (234), (243)

(12)(34), (13)(24), (14)(23)

(1234), (1243), (1324), (1342), (1423), (1432)

Version: 1 Owner: rmilson Author(s): rmilson

307.2 permutation group

**A permutation group is a pair (G, X) where G is an abstract group, and X is a set on
**

which G acts faithfully. Alternatively, this can be thought of as a group G equipped with a

homomorphism in to Sym(X), the symmetric group on X.

Version: 2 Owner: bwebste Author(s): bwebste

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Chapter 308

20B15 – Primitive groups

308.1 primitive transitive permutation group

1: A finite set

2: G transitive permutation group on A

3: ∀B ⊂ A block or B = 1

example

1: S4 is a primitive transitive permutation group on {1, 2, 3, 4}

counterexample

1: D8 is not a primitive transitive permutation group on the vertices of a square

stabilizer maximal necessary and sufficient for primitivity

1: A finite set

2: G transitive permutation group on A

3: G primitive ⇔ ∀a ∈ A : H 6 G & H ⊃ StabG (a) ⇒ H = G or H = StabG (a)

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Note: This was a “seed” entry written using a short-hand format described in this FAQ.

Version: 4 Owner: Thomas Heye Author(s): yark, apmxi

1285

Chapter 309

**20B20 – Multiply transitive finite
**

groups

309.1 Jordan’s theorem (multiply transitive groups)

Let G be a sharply n-transitive permutation group, with n 6 4. Then

1. G is similar to Sn with the standard action or

2. n = 4 and G is similar to M11 , the Mathieu group of degree 10 or

3. n = 5 and G is similar to M12 , the Mathieu group of degree 11.

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309.2 multiply transitive

**Let G be a group, X a set on which it acts. Let X (n) be the set of order n-tuples of distinct
**

elements of X. This is a G-set by the diagonal action:

g · (x1 , . . . , xn ) = (g · x1 , . . . , g · xn )

The action of G on X is said to be n-transitive if it acts transitively on X (n) .

**For example, the standard action of S n , the symmetric group, is n-transitive, and the stan-
**

dard action of An , the alternating group, is (n − 2)-transitive.

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309.3 sharply multiply transitive

**Let G be a group, and X a set that G acts on, and let X (n) be the set of order n-tuples of
**

distinct elements of X. Then the action of G on X is sharply n-transitive if G acts regularly

on X (n) .

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1287

Chapter 310

**20B25 – Finite automorphism groups
**

of algebraic, geometric, or

combinatorial structures

310.1 diamond theory

**Diamond theory is the theory of affine groups over GF (2) acting on small square and cubic
**

arrays. In the simplest case, the symmetric group of order 4 acts on a two-colored Diamond

figure like that in Plato’s Meno dialogue, yielding 24 distinct patterns, each of which has

some ordinary or color-interchange symmetry.

**This can be generalized to (at least) a group of order approximately 1.3 trillion acting on a
**

4x4x4 array of cubes, with each of the resulting patterns still having nontrivial symmetry.

**The theory has applications to finite geometry and to the construction of the large Witt
**

design underlying the Mathieu group of degree 24.

Further Reading

• ”Diamond Theory,” http://m759.freeservers.com/

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1288

Chapter 311

20B30 – Symmetric groups

311.1 symmetric group

**Let X be a set. Let S(X) be the set of permutations of X (i.e. the set of bijective functions
**

on X). Then the act of taking the composition of two permutations induces a group structure

on S(X). We call this group the symmetric group and it is often denoted Sym(X).

Version: 5 Owner: bwebste Author(s): bwebste, antizeus

311.2 symmetric group

**Let X be a set. Let S(X) be the set of permutations of X (i.e. the set of bijective functions
**

on X). Then the act of taking the composition of two permutations induces a group structure

on S(X). We call this group the symmetric group and it is often denoted Sym(X).

**When X has a finite number n of elements, we often refer to the symmetric group as Sn ,
**

and describe the elements by using cycle notation.

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Chapter 312

**20B35 – Subgroups of symmetric
**

groups

312.1 Cayley’s theorem

Let G be a group, then G is isomorphic to a subgroup of the permutation group SG

**If G is finite and of order n, then G is isomorphic to a subgroup of the permutation group
**

Sn

**Furthermore, suppose H is a proper subgroup of G. Let X = {Hg|g ∈ G} be the set of
**

right cosets in G. The map θ : G → SX given by θ(x)(Hg) = Hgx is a homomorphism. The

kernel is the largest normal subgroup of H. We note that |SX | = [G : H]!. Consequently if

|G| doesn’t divide [G : H]! then θ is not an isomorphism so H contains a non-trivial normal

subgroup, namely the kernel of θ.

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Chapter 313

20B99 – Miscellaneous

313.1 (p, q) shuffle

Definition.

Let p and q be positive natural numbers. Further, let S(k) be the set of permutations of the

numbers {1, . . . , k}. A permutation τ ∈ S(p + q) is a (p, q) shuffle if

τ (1) < · · · < τ (p),

τ (p + 1) < · · · < τ (p + q).

The set of all (p, q) shuffles is denoted by S(p, q).

**It is clear that S(p, q) ⊂ S(p + q). Since a (p, q) shuffle is completely
**

determined by how

p+q

the p first elements are mapped, the cardinality of S(p, q) is p . The wedge product of a

p-form and a q-form can be defined as a sum over (p, q) shuffles.

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313.2 Frobenius group

**A permutation group G on a set X is Frobenius if no non-trivial element of G fixes more
**

than one element of X. Generally, one also makes the restriction that at least one non-trivial

element fix a point. In this case the Frobenius group is called non-regular.

**The stabilizer of any point in X is called a Frobenius complement, and has the remarkable
**

property that it is distinct from any conjugate by an element not in the subgroup. Conversely,

if any finite group G has such a subgroup, then the action on cosets of that subgroup makes

G into a Frobenius group.

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313.3 permutation

**A permutation of a set {a1 , a2 , . . . , an } is an arrangement of its elements. For example, if
**

S = {ABC} then ABC, CAB , CBA are three different permutations of S.

The number of permutations of a set with n elements is n!.

**A permutation can also be seen as a bijective function of a set into itself. For example, the
**

permutation CAB could be seen a function that assigns:

f (A) = C, f (B) = A, f (C) = B.

**In fact, every bijection of a set into itself gives a permutation, and any permutation gives
**

rise to a bijective function.

**Therefore, we can say that there are n! bijective fucntion from a set with n elements into
**

itself.

**Using the function approach, it can be proved that any permutation can be expressed
**

as a composition of disjoint cycles and also as composition of (not necessarily disjoint)

transpositions.

**Moreover, if σ = τ1 τ2 · · · τm = ρ1 ρ2 · · · ρn are two factorization of a permutation σ into
**

transpositions, then m and n must be both even or both odd. So we can label permutations

as even or odd depending on the number of transpositions for any decomposition.

**Permutations (as functions) form a non-abelian group with function composition as binary operation
**

called symmetric group of order n. The subset of even permutations becomes a subgroup

called the alternating group of order n.

Version: 3 Owner: drini Author(s): drini

313.4 proof of Cayley’s theorem

**Let G be a group, and let SG be the permutation group of the underlying set G. For each
**

g ∈ G, define ρg : G → G by ρg (h) = gh. Then ρg is invertible with inverse ρg−1 , and so is a

permutation of the set G.

Define Φ : G → SG by Φ(g) = ρg . Then Φ is a homomorphism, since

(Φ(gh))(x) = ρgh (x) = ghx = ρg (hx) = (ρg ◦ ρh )(x) = ((Φ(g))(Φ(h)))(x)

**And Φ is injective, since if Φ(g) = Φ(h) then ρg = ρh , so gx = hx for all x ∈ X, and so
**

g = h as required.

1292

So Φ is an embedding of G into its own permutation group. If G is finite of order n, then

simply numbering the elements of G gives an embedding from G to Sn .

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1293

Chapter 314

**20C05 – Group rings of finite groups
**

and their modules

314.1 group ring

**For any group G, the group ring Z[G] is defined to be the ring whose additive group is the
**

abelian group of formal integer linear combinations of elements of G, and whose multiplica-

tion operation is defined by multiplication in G, extended Z–linearly to Z[G].

**More generally, for any ring R, the group ring of G over R is the ring R[G] whose additive
**

group is the abelian group of formal R–linear combinations of elements of G, i.e.:

( n )

X

R[G] := ri gi ri ∈ R, gi ∈ G ,

i=1

**and whose multiplication operation is defined by R–linearly extending the group multiplica-
**

tion operation of G. In the case where K is a field, the group ring K[G] is usually called a

group algebra.

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1294

Chapter 315

**20C15 – Ordinary representations and
**

characters

315.1 Maschke’s theorem

**Let G be a finite group, and k a field of characteristic not dividing |G|. Then any representation
**

V of G over k is completely reducible.

**W e need only show that any subrepresentation has a compliment, and the result follows
**

by induction.

**Let V be a representation of G and W a subrepresentation. Let π : V → W be an arbitrary
**

projection, and let

1 X −1

π 0 (v) = g π(gv)

|G| g∈G

**This map is obviously G-equivariant, and is the identity on W , and its image is contained in
**

W , since W is invariant under G. Thus it is an equivariant projection to W , and its kernel

is a compliment to W .

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315.2 a representation which is not completely reducible

**If G is a finite group, and k is a field whose characteristic does divide the order of the group,
**

then Maschke’s theorem fails. For example let V be the regular representation of G, which

can be thought of as functions from G to k, with the G action g · ϕ(g 0) = ϕ(g −1 g 0). Then

this representation is not completely reducible.

1295

There is an obvious trivial subrepresentation W of V , consisting of the constant functions. I

claim that there is no complementary invariant subspace to this one. If W 0 is such a subspace,

then there is a homomorphism ϕ : V → V /W 0 ∼ = k. Now consider the characteristic function

of the identity e ∈ G (

1 g=e

δe (g) =

0 g 6= e

and ` = ϕ(δe ) in V /W 0. This is not zero since δ generates the representation V . By G-

equivarience, ϕ(δg ) = ` for all g ∈ G. Since

X

η= η(g)δg

g∈G

for all η ∈ V , !

X

0

W = ϕ(η) = ` η(g) .

g∈G

Thus, X

ker ϕ = {η ∈ V | η(g) = 0}.

∈G

**But since the characteristic of the field k divides the order of G, W 6 W 0 , and thus could
**

not possibly be complimentary to it.

**For example, if G = C2 = {e, f } then the invariant subspace of V is spanned by e + f . For
**

characteristics other than 2, e − f spans a complimentary subspace, but over characteristic

2, these elements are the same.

Version: 1 Owner: bwebste Author(s): bwebste

315.3 orthogonality relations

**First orthogonality relations: Let χ1 , χ2 be characters of representations V1 , V2 of a finite group
**

G over a field k of characteristic 0. Then

1 X

(χ1 , χ2 ) = χ1 (g)χ2 (g) = dim(HomV1 V2 ).

|G| g∈G

**F irst of all, consider the special case where V = k with the trivial action of the group.
**

Then HomG (k, V2 ) ∼= V2G , the fixed points. On the other hand, consider the map

1 X

φ= g : V2 → V2

|G| g∈G

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(with the sum in End(V2 )). Clearly, the image of this map is contained in V2G , and it is

the identity restricted to V2G . Thus, it is a projection with image V2G . Now, the rank of a

projection (over a field of characteristic 0) is its trace. Thus,

1 X

dimk HomG (k, V2 ) = dim V2G = tr(φ) = χ2 (g)

|G|

which is exactly the orthogonality formula for V1 = k.

**Now, in general, Hom(V1 , V2 ) ∼
**

= V1∗ ⊗V2 is a representation, and HomG (V1 , v2 ) = (Hom(V1 , V2 ))G .

Since χV1∗ ⊗V2 = χ1 χ2 ,

X

dimk HomG (V1 , V2 ) = dimk (Hom(V1 , V2 ))G = χ1 χ2

g∈G

which is exactly the relation we desired.

**In particular, if V1 , V2 irreducible, by Schur’s lemma
**

(

D V1 ∼

= V2

HomV1 V2 =

0 V1 V2

**where D is a division algebra. In particular, non-isomorphic irreducible representations have
**

orthogonal characters. Thus, for any representation V , the multiplicities ni in the unique

decomposition of V into the direct sum of irreducibles

V ∼

= V1⊕n1 ⊕ · · · ⊕ Vm⊕nm

**where Vi ranges over irreducible representations of G over k, can be determined in terms of
**

the character inner product:

(ψ, χi )

ni =

(χi , χi )

**where ψ is the character of V and χi the character of Vi . In particular, representations over
**

a field of characteristic zero are determined by their character. Note: This is not true over

fields of positive characteristic.

**If the field k is algebraically closed, the only finite division algebra over k is k itself, so the
**

characters of irreducible representations form an orthonormal basis for the vector space of

class functions with respect to this inner product. Since (χi , χi ) = 1 for all irreducibles, the

multiplicity formula above reduces to ni = (ψ, χi ).

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Second orthogonality relations: We assume now that k is algebraically closed. Let g, g 0 be

elements of a finite group G. Then

(

X |CG (g1 )| g ∼ g 0

χ(g)χ(g 0) =

χ

0 g g0

**where the sum is over the characters of irreducible representations, and CG (g) is the centralizer
**

of g.

**L et χ1 , . . . , χn be the characters of the irreducible representations, and let g1 , . . . , gn be
**

representatives of the conjugacy classes.

p

Let A be the matrix whose ijth entry is |G : CG (gj )|(χi (gj )). By first orthogonality,

AA∗ = |G|I (here ∗ denotes conjugate transpose), where I is the identity matrix. Since left

inverses are right inverses, A∗ A = |G|I. Thus,

n

X

p

|G : CG (gi )||G : CG (gk )| χj (gi )χj (gk ) = |G|δij .

j=1

**Replacing gi or gk with any conjuagate will not P change the expression above. thus, if our
**

two elements are not conjugate, we obtain that χ χ(g)χ(g 0) = 0. On the other hand, if

g ∼ g 0 , then i = k in the sum above, which reduced to the expression we desired.

P

A special case of this result, applied to 1 is that |G| = χ χ(1)2 , that is, the sum of the

squares of the dimensions of the irreducible representations of any finite group is the order

of the group.

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1298

Chapter 316

**20C30 – Representations of finite
**

symmetric groups

316.1 example of immanent

If χ = 1 we obtain the permanent. If χ = sgn we obtain the determinant.

Version: 1 Owner: gholmes74 Author(s): gholmes74

316.2 immanent

**Let χ : Sn → C be a complex character. For any n × n matrix A define
**

X n

Y

Immχ (A) = χ(σ) A(j, σj)

σ∈Sn j=1

functions obtained in this way are called immanents.

Version: 4 Owner: gholmes74 Author(s): gholmes74

316.3 permanent

**The permanent of an n × n matrix A over C is the number
**

n

XY

per(A) = A(j, σj)

σ∈Sn j=1

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1300

Chapter 317

20C99 – Miscellaneous

317.1 Frobenius reciprocity

**Let V be a finite-dimensional representation of a finite group G, and let W be a represen-
**

tation of a subgroup H ⊂ G. Then the characters of V and W satisfy the inner product

relation

(χInd(W ) , χV ) = (χW , χRes(V ) )

where Ind and Res denote the induced representation IndG

H and the restriction representation

ResG

H .

**The Frobenius reciprocity theorem is often given in the stronger form which states that
**

Res and Ind are adjoint functors between the category of G–modules and the category of

H–modules:

HomH (W, Res(V )) = HomG (Ind(W ), V ),

or, equivalently

V ⊗ Ind(W ) = Ind(Res(V ) ⊗ W ).

Version: 4 Owner: djao Author(s): rmilson, djao

317.2 Schur’s lemma

**Schur’s lemma in representation theory is an almost trivial observation for irreducible modules,
**

but deserves respect because of its profound applications and implications.

**Lemma 5 (Schur’s lemma). Let G be a finite group represented on irreducible G-modules
**

V and W . Any G-module homomorphism f : V → W is either invertible or the zero map.

1301

T he only insight here is that both ker f and im f are G-submodules of V and W respec-

tively. This is routine. However, because V is irreducible, ker f is either trivial or all of V .

In the former case, im f is all of W , also because W is irreducible, so f is invertible. In the

latter case, f is the zero map.

**The following corollary is a very useful form of Schur’s lemma, in case that our representations
**

are over an algebraically closed field.

Corollary 1. If G is represented over an algebraically closed field F on irreducible G-modules

V and W , then any G-module homomorphism f : V → W is a scalar.

**T he insight in this case is to consider the modules V and W as vector spaces over F . Notice
**

then that the homomorphism f is a linear transformation and therefore has an eigenvalue λ

in our algebraically closed F . Hence, f −λ1 is not invertible. By Schur’s lemma, f −λ1 = 0.

In other words, f = λ, a scalar.

Version: 14 Owner: rmilson Author(s): rmilson, NeuRet

317.3 character

**Let ρ : G −→ GL(V ) be a finite dimensional representation of a group G (i.e., V is a
**

finite dimensional vector space over its scalar field K). The character of ρ is the function

χV : G −→ K defined by

χV (g) := Tr(ρ(g))

where Tr is the trace function.

Properties:

**• χV (g) = χV (h) if g is conjugate to h in G. (Equivalently, a character is a class function
**

on G.)

**• If G is finite, the characters of the irreducible representations of G over the complex numbers
**

form a basis of the vector space of all class functions on G (with pointwise addition

and scalar multiplication).

**• Over the complex numbers, the characters of the irreducible representations of G are
**

orthonormal under the inner product

1 X

(χ1 , χ2 ) := χ1 (g)χ2 (g)

|G| g∈G

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1302

317.4 group representation

**Let G be a group, and let V be a vector space. A representation of G in V is a group homomorphism
**

ρ : G −→ GL(V ) from G to the general linear group GL(V ) of invertible linear transformations

of V .

**Equivalently, a representation of G is a vector space V which is a (left) module over the
**

group ring Z[G]. The equivalence is achieved by assigning to each homomorphism ρ : G −→

GL(V ) the module structure whose scalar multiplication is defined by g · v := (ρ(g))(v), and

extending linearly.

Special kinds of representations (preserving all notation from above)

A representation is faithful if either of the following equivalent conditions is satisfied:

• ρ : G −→ GL(V ) is injective

• V is a faithful left Z[G]–module

**A subrepresentation of V is a subspace W of V which is a left Z[G]–submodule of V ; or,
**

equivalently, a subspace W of V with the property that

(ρ(g))(w) ∈ W for all w ∈ W.

**A representation V is called irreducible if it has no subrepresentations other than itself and
**

the zero module.

Version: 2 Owner: djao Author(s): djao

317.5 induced representation

**Let G be a group, H ⊂ G a subgroup, and V a representation of H, considered as a Z[H]–
**

module. The induced representation of ρ on G, denoted IndG

H (V ), is the Z[G]–module whose

underlying vector space is the direct sum

M

σV

σ∈G/H

**of formal translates of V by left cosets σ in G/H, and whose multiplication operation is
**

defined by choosing a set {gσ }σ∈G/H of coset representatives and setting

g(σv) := τ (hv)

1303

where τ is the unique left coset of G/H containing g · gσ (i.e., such that g · gσ = gτ · h for

some h ∈ H).

**One easily verifies that the representation IndG
**

H (V ) is independent of the choice of coset

representatives {gσ }.

Version: 1 Owner: djao Author(s): djao

317.6 regular representation

**Given a group G, the regular representation of G over a field K is the representation
**

ρ : G −→ GL( K G ) whose underlying vector space K G is the K–vector space of formal

linear combinations of elements of G, defined by

n

! n

X X

ρ(g) kigi := ki (ggi)

i=1 i=1

for ki ∈ K, g, gi ∈ G.

**Equivalently, the regular representation is the induced representation on G of the trivial
**

representation on the subgroup {1} of G.

Version: 2 Owner: djao Author(s): djao

317.7 restriction representation

**Let ρ : G −→ GL( V ) be a representation on a group G. The restriction representation of ρ
**

to a subgroup H of G, denoted ResG H (V ), is the representation ρ|H : H −→ GL( V ) obtained

by restricting the function ρ to the subset H ⊂ G.

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1304

Chapter 318

**20D05 – Classification of simple and
**

nonsolvable groups

318.1 Burnside p − q theorem

**If a finite group G is not solvable, the order of G is divisible by at least 3 distinct primes.
**

Alternatively, any groups whose order is divisible by only two distinct primes is solvable

(these two distinct primes are the p and q of the title).

Version: 2 Owner: bwebste Author(s): bwebste

318.2 classification of semisimple groups

**For every semisimple group G there is a normal subgroup H of G, (called the centerless com-
**

petely reducible radical) which isomorphic to a direct product of nonabelian simple groups

such that conjugation on H gives an injection into Aut(H). Thus G is isomorphic to a

subgroup of Aut(H) containing the inner automorphisms, and for every group H isomorphic

to a direct product of non-abelian simple groups, every such subgroup is semisimple.

Version: 1 Owner: bwebste Author(s): bwebste

318.3 semisimple group

**A group G is called semisimple if it has no proper normal solvable subgroups. Every group
**

is an extension of a semisimple group by a solvable one.

1305

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1306

Chapter 319

**20D08 – Simple groups: sporadic
**

groups

319.1 Janko groups

**The Janko groups denoted by J1 , J2 , J3 , and J4 are four of the 26 sporadic groups. They were
**

discovered by Z. Janko in 1966 and published in the article ”A new finite simple group with

abelan Sylow subgroups and its characterization.” (Journal of algebra, 1966, 32: 147-186).

Each of these groups have very intricate matrix representations as maps into large general linear groups.

For example, the matrix K corresponding to J4 gives a representation of J4 in GL112 (2).

Version: 7 Owner: mathcam Author(s): mathcam, Thomas Heye

1307

Chapter 320

**20D10 – Solvable groups, theory of
**

formations, Schunck classes, Fitting

classes, π-length, ranks

320.1 Čuhinin’s Theorem

**Let G be a finite, π-separable group, for some set π of primes. Then if H is a maximal
**

π-subgroup of G, the index of H in G, |G : H|, is coprime to all elements of π and all such

subgroups are conjugate. Such a subgroup is called a Hall π-subgroup. For π = {p}, this

essentially reduces to the Sylow theorems (with unnecessary hypotheses).

**If G is solvable, it is π-separable for all π, so such subgroups exist for all π. This result is
**

often called Hall’s theorem.

Version: 4 Owner: bwebste Author(s): bwebste

320.2 separable

Let π be a set of primes. A finite group G is called π-separable if there exists a composition series

{1} = G0 · · · Gn = G

**such that Gi+1 /Gi is a π-group, or a π 0 -group. π-separability can be thought of as a gener-
**

alization of solvability; a group is π-separable for all sets of primes if and only it is solvable.

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320.3 supersolvable group

A group G is supersolvable if it has a finite normal series

G = G0 G1 · · · Gn = 1

with the property that each factor group Gi−1 /Gi is cyclic.

A supersolvable group is solvable.

Finitely generated nilpotent groups are supersolvable.

Version: 1 Owner: mclase Author(s): mclase

1309

Chapter 321

20D15 – Nilpotent groups, p-groups

321.1 Burnside basis theorem

**If G is a p-group, then Frat G = G0 Gp , where Frat G is the Frattini subgroup, G0 the
**

commutator subgroup, and Gp is the subgroup generated by p-th powers.

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1310

Chapter 322

**20D20 – Sylow subgroups, Sylow
**

properties, π-groups, π-structure

322.1 π-groups and π 0 -groups

**Let π be a set of primes. A finite group G is called a π-group if all the primes dividing |G|
**

are elements of π, and a π 0 -group if none of them are. Typically, if π is a singleton π = {p},

we write p-group and p0 -group for these.

Version: 2 Owner: bwebste Author(s): bwebste

322.2 p-subgroup

**Let G be a finite group with order n, and let p be a prime integer. We can write n = pk m
**

for some k, m integers, such that k and m are coprimes (that is, pk is the highest power of p

that divides n). Any subgroup of G whose order is pk is called a Sylow p-subgroup or simply

p-subgroup.

While there is no reason for p-subgroups to exist for any finite group, the fact is that all

groups have p-subgroups for every prime p that divides |G|. This statement is the First

Sylow theorem When |G| = pk we simply say that G is a p-group.

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1311

322.3 Burnside normal complement theorem

**Let G be a finite group, and S a Sylow subgroup such that CG (S) = NG (S). Then
**

T S has a

normal complement. That is, there exists a normal subgroup N G such that S N = {1}

and SN = G.

Version: 1 Owner: bwebste Author(s): bwebste

322.4 Frattini argument

**If H is a normal subgroup of a finite group G, and S is a Sylow subgroup of H, then G =
**

HNG (S), where NG (S) is the normalizer of S in G.

Version: 1 Owner: bwebste Author(s): bwebste

322.5 Sylow p-subgroup

**If (G, ∗) is a group then any subgroup of order pa for any integer a is called a p-subgroup.
**

If|G| = pa m, where p - m then any subgroup S of G with |S| = pa is a Sylow p-subgroup.

We use Sylp (G) for the set of Sylow p-groups of G.

Version: 3 Owner: Henry Author(s): Henry

322.6 Sylow theorems

**Let G be a finite group whose order is divisible by the prime p. Suppose pm is the highest
**

power of p which is a factor of |G| and set k = p|G|

m

• The group G contains at least one subgroup of order pm

• Any two subgroups of G of order pm are conjugate

**• The number of subgroups of G of order pm is congruent to 1 modulo p and is a factor
**

of k

Version: 1 Owner: vitriol Author(s): vitriol

1312

322.7 Sylow’s first theorem

existence of subgroups of prime-power order

1: G finite group

2: p prime

3: pk divides |G|

4: ∃H 6 G : |H| = pk

Note: This is a “seed” entry written using a short-hand format described in this FAQ.

Version: 2 Owner: bwebste Author(s): yark, apmxi

322.8 Sylow’s third theorem

**Let G finite group, and let n be the number of Sylow p-subgroups of G. Then n ⇔ 1 (mod p),
**

and any two Sylow p-subgroups of G are conjugate to one another.

Version: 8 Owner: bwebste Author(s): yark, apmxi

**322.9 application of Sylow’s theorems to groups of or-
**

der pq

**We can use Sylow’s theorems to examine a group G of order pq, where p and q are primes
**

and p < q.

**Let nq denote the number of Sylow q-subgroups of G. Then Sylow’s theorems tell us that nq
**

is of the form 1 + kq for some integer k and nq divides pq. But p and q are prime and p < q,

so this implies that nq = 1. So there is exactly one Sylow q-subgroup, which is therefore

normal (indeed, characteristic) in G.

T

Denoting the Sylow q-subgroup by Q, and letting P be a Sylow p-subgroup, then Q P =

{1} and QP = G, so G is a semidirect product of Q and P . In particular, if there is only

one Sylow p-subgroup, then G is a direct product of Q and P , and is therefore cyclic.

Version: 9 Owner: yark Author(s): yark, Manoj, Henry

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322.10 p-primary component

**Definition 27. Let G be a finite abelian group and let p ∈ N be a prime. The p-primary
**

component of G, Πp , is the subgroup of all elements whose order is a power of p.

**Note: The p-primary component of an abelian group G coincides with the unique Sylow
**

p-subgroup of G.

Version: 2 Owner: alozano Author(s): alozano

322.11 proof of Frattini argument

**Let g ∈ G be any element. Since H is normal, gSg −1 ⊂ H. Since S is a Sylow subgroup
**

of H, gSg −1 = hSh−1 for some h ∈ H, by Sylow’s theorems. Thus n = h−1 g normalizes S,

and so g = hn for h ∈ H and n ∈ NG (S).

Version: 1 Owner: bwebste Author(s): bwebste

322.12 proof of Sylow theorems

We let G be a group of order pm k where p - k and prove Sylow’s theorems.

**First, a fact which will be used several times in the proof:
**

Proposition 8. If p divides the size of every conjugacy class outside the center then p divides

the order of the center.

**Proof: This follows from this Centralizer:
**

X

|G| = Z(G) + |[a]|

a∈Z(G)

/

If p divides the left hand side, and divides all but one entry on the right hand side, it must

divide every entry on the right side of the equation, so p|Z(G).

Proposition 9. G has a Sylow p-subgroup

**Proof: By induction on |G|. If |G| = 1 then there is no p which divides its order, so the
**

condition is trivial.

**Suppose |G| = pm k, p - k, and the proposition holds for all groups of smaller order. Then
**

we can consider whether p divides the order of the center, Z(G).

1314

If it does then, by Cauchy’s theorem, there is an element of Z(G) of order p, and therefore

a cyclic subgroup generated by p, hpi, also of order p. Since this is a subgroup of the center,

it is normal, so G/hpi is well-defined and of order pm−1 k. By the inductive hypothesis, this

group has a subgroup P/hpi of order pm−1 . Then there is a corresponding subgroup P of G

which has |P | = |P/hpi| · |N| = pm .

**On the other hand, if p - |Z(G)| then consider the conjugacy classes not in the center. By
**

the proposition above, since Z(G) is not divisible by p, at least one conjugacy class can’t

be. If a is a representative of this class then we have p - |[a]| = [G : C(a)], and since

|C(a)| · [G : C(a)] = |G|, pm | |C(a)|. But C(a) 6= G, since a ∈

/ Z(G), so C(a) has a subgroup

m

of order p , and this is also a subgroup of G.

**Proposition 10. The intersection of a Sylow p-subgroupTwith the normalizer
**

T of a Sylow

p-subgroup is the intersection of the subgroups. That is, Q NG (P ) = Q P .

T T

Proof: If P and Q are Sylow p-subgroups, consider R = Q NG (P ). Obviously Q P ⊆ R.

In addition, since R ⊆ NG (P ), the second isomorphism theorem tells us that RP is a group,

|R|·|P

T | . P is a subgroup of RP , so pm | |RP |. But |R| is a subgroup of Q and P

and |RP | = |R P|

m

is a Sylow p-subgroup, so |R| · |P | is a multiple of p. Then it must

T be that |RP | = p , and

therefore P = RP , and so R ⊆ P . Obviously R ⊆ Q, so R ⊆ Q P .

The following construction will be used in the remainder of the proof:

**Given any Sylow p-subgroup P , consider the set of its conjugates C. Then X ∈ C ↔ X =
**

xP x−1 = {xpx−1 |∀p ∈ P } for some x ∈ G. Observe that every X ∈ C is a Sylow p-subgroup

(and we will show that the converse holds as well). We define a group action of a subset G

on C by:

g · X = g · xP x−1 = gxP x−1 g −1 = (gx)P (gx)−1

This is clearly a group action, so we can consider the orbits of P under it. Of course, if all G

is used then there is only one orbit, so we restrict the action to a Sylow p-subgroup Q. Name

the orbits O1 , . . . , Os , and let P1 , . . . , Ps be representatives of the corresponding orbits. By

the orbit-stabilizer theorem, the size of an orbit isTthe index of the T stabilizer, and under

T this

action the stabilizer of any Pi is just NQ (Pi ) = Q NG (Pi ) = Q P , so |Oi | = [Q : Q Pi ].

T

There are two easy T results on this construction. If Q = Pi then |Oi| = [Pi : Pi Pi ] = 1. If

Q 6= Pi then [Q : Q Pi ] > 1, and since the index of any subgroup of Q divides Q, p | |Oi|.

**Proposition 11. The number of conjugates of any Sylow p-subgroup of G is congruent to 1
**

modulo p

**In the construction above, let Q = P1 . Then |O1 | = 1 and p | |Oi| for i 6= 1. Since the
**

number of conjugates of P is the sum of the number in each orbit, the number of conjugates

is of the form 1 + k2 p + k3 p + · · · + ks p, which is obviously congruent to 1 modulo p.

Proposition 12. Any two Sylow p-subgroups are conjugate

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Proof: Given a Sylow p-subgroup P and any other Sylow p-subgroup Q, consider again

the construction given above. If Q is not conjugate to P then Q 6= Pi for every i, and

therefore p | |Oi| for every orbit. But then the number of conjugates of P is divisible by p,

contradicting the previous result. Therefore Q must be conjugate to P .

**Proposition 13. The number of subgroups of G of order pm is congruent to 1 modulo p and
**

is a factor of k

**Proof: Since conjugates of a Sylow p-subgroup are precisely the Sylow p-subgroups, and since
**

a Sylow p-subgroup has 1 modulo p conjugates, there are 1 modulo p Sylow p-subgroups.

**Since the number of conjugates is the index of the normalizer, it must be |G : NG (P )|. Since
**

P is a subgroup of its normalizer, pm | NG (P ), and therefore |G : NG (P )| | k.

Version: 3 Owner: Henry Author(s): Henry

**322.13 subgroups containing the normalizers of Sylow
**

subgroups normalize themselves

**Let G be a finite group, and S a Sylow subgroup. Let M be a subgroup such that NG (S) ⊂
**

M. Then M = NG (M).

**B y order considerations, S is a Sylow subgroup of M. Since M is normal in NG (M), by
**

the Frattini argument, NG (M) = NG (S)M = M.

Version: 3 Owner: bwebste Author(s): bwebste

1316

Chapter 323

**20D25 – Special subgroups (Frattini,
**

Fitting, etc.)

323.1 Fitting’s theorem

**If G is a finite group and M and N are normal nilpotent subgroups, then MN is also a
**

normal nilpotent subgroup.

Thus, any finite group has a maximal normal nilpotent subgroup, called its Fitting subgroup.

Version: 1 Owner: bwebste Author(s): bwebste

323.2 characteristically simple group

**A group G is called characterisitically simple if its only characteristic subgroups are {1}
**

and G. Any finite characteristically simple group is the direct product of several copies of

isomorphic simple groups.

Version: 3 Owner: bwebste Author(s): bwebste

323.3 the Frattini subgroup is nilpotent

The Frattini subgroup Frat G of any finite group G is nilpotent.

**L et S be a Sylow p-subgroup of G. Then by the Frattini argument, (Frat G)NG (S) = G.
**

Since the Frattini subgroup is formed of non-generators, NG (S) = G. Thus S is normal in

1317

G, and thus in Frat G. Any subgroup whose Sylow subgroups are all normal is nilpotent.

Version: 4 Owner: bwebste Author(s): bwebste

1318

Chapter 324

**20D30 – Series and lattices of
**

subgroups

324.1 maximal condition

**A group is said to satisfy the maximal condition if every strictly ascending chain of
**

subgroups

G1 ⊂ G2 ⊂ G3 ⊂ · · ·

is finite.

This is also called the ascending chain condition.

A group satifies the maximal condition if and only if the group and all its subgroups are

finitely generated.

**Similar properties are useful in other classes of algebraic structures: see for example the
**

noetherian condition for rings and modules.

Version: 2 Owner: mclase Author(s): mclase

324.2 minimal condition

**A group is said to satisfy the minimal condition if every strictly descending chain of
**

subgroups

G1 ⊃ G2 ⊃ G3 ⊃ · · ·

is finite.

This is also called the descending chain condition.

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A group which satisfies the minimal condition is necessarily periodic. For if it contained an

element x of infinite order, then

n

hxi ⊃ hx2 i ⊃ hx4 i ⊃ · · · ⊃ hx2 i ⊃ · · ·

is an infinite descending chain of subgroups.

**Similar properties are useful in other classes of algebraic structures: see for example the
**

artinian condition for rings and modules.

Version: 1 Owner: mclase Author(s): mclase

324.3 subnormal series

Let G be a group with a subgroup H, and let

G = G0 G1 · · · Gn = H (324.3.1)

**be a series of subgroups with each Gi a normal subgroup of Gi−1 . Such a series is called a
**

subnormal series or a subinvariant series.

If in addition, each Gi is a normal subgroup of G, then the series is called a normal series.

**A subnormal series in which each Gi is a maximal normal subgroup of Gi−1 is called a
**

composition series.

**A normal series in which Gi is a maximal normal subgroup of G contained in Gi−1 is called
**

a principal series or a chief series.

**Note that a composition series need not end in the trivial group 1. One speaks of a composi-
**

tion series (1) as a composition series from G to H. But the term composition series

for G generally means a compostion series from G to 1.

Similar remarks apply to principal series.

Version: 1 Owner: mclase Author(s): mclase

1320

Chapter 325

20D35 – Subnormal subgroups

325.1 subnormal subgroup

Let G be a group, and H a subgroup of G. Then H is subnormal if there exists a finite series

H = H0 hdH1 hd · · · hdtHn = G

with Hi a normal subgroup of Hi+1 .

Version: 1 Owner: bwebste Author(s): bwebste

1321

Chapter 326

20D99 – Miscellaneous

326.1 Cauchy’s theorem

**Let G be a finite group and let p be a prime dividing |G|. Then there is an element of G of
**

order p.

Version: 1 Owner: Evandar Author(s): Evandar

326.2 Lagrange’s theorem

**Let G be a finite group and let H be a subgroup of G. Then the order of H divides the
**

order of G.

Version: 2 Owner: Evandar Author(s): Evandar

326.3 exponent

**If G is a finite group, then the exponent of G, denoted exp G, is the smallest positive integer
**

n such that, for every g ∈ G, g n = eG . Thus, for every group G, exp G divides G, and, for

every g ∈ G, |g| divides exp G.

The concept of exponent for finite groups is similar to that of characterisic for rings.

**If G is a finite abelian group, then there exists g ∈ G with |g| = exp G. As a result of the
**

fundamental theorem of finite abelian groups, there exist a1 , . . . , an with ai dividing ai+1 for

every integer i between 1 and n such that G ∼ = Za1 ⊕ · · · ⊕ Zan . Since, for every c ∈ G,

1322

can = eG , then exp G ≤ an . Since |(0, . . . , 0, 1)| = an , then exp G = an , and the result

follows.

Following are some examples of exponents of nonabelian groups.

Since |(12)| = 2, |(123)| = 3, and |S3 | = 6, then exp S3 = 6.

**In Q8 = {1, −1, i, −i, j, −j, k, −k}, the ring of quaternions of order eight, since |i| = | − i| =
**

|j| = | − j| = |k| = | − k| = 4 and 14 = (−1)4 = 1, then exp Q = 4.

**Since the order of a product of two disjoint transpositions is 2, the order of a three cycle is
**

3, and the only nonidentity elements of A4 are products of two disjoint transpositions and

three cycles, then exp A4 = 6.

**Since |(123)| = 3 and |(1234)| = 4, then exp S4 ≥ 12. Since S4 is not abelian, then it is not
**

cyclic, and thus contains no element of order 24. It follows that exp S4 = 12.

Version: 5 Owner: Wkbj79 Author(s): Wkbj79

326.4 fully invariant subgroup

A subgroup H of a group G is fully invariant if f (H) ⊆ H for all endomorphisms f : G → G

This is a stronger condition than being a characteristic subgroup.

The derived subgroup is fully invariant.

Version: 1 Owner: mclase Author(s): mclase

326.5 proof of Cauchy’s theorem

**Let G be a finite group and p be a prime divisor of |G|. Consider the set X of all ordered
**

strings (x1 , x2 , . . . , xp ) for which x1 x2 . . . xp = e. Note |X| = |G|p−1, i.e. a multiple of

p. There is a natural group action of Zp on X. m ∈ Zp sends the string (x1 , x2 , . . . , xp )

to (xm+1 , . . . , xp , x1 , . . . , xm ). By orbit-stabilizer theorem each orbit contains exactly 1 or

p strings. Since (e, e, . . . , e) has an orbit of cardinality 1, and the orbits partition X, the

cardinality of which is divisible by p, there must exist at least one other string (x1 , x2 , . . . , xp )

which is left fixed by every element of Zp . i.e. x1 = x2 = . . . = xp and so there exists an

element of order p as required.

Version: 1 Owner: vitriol Author(s): vitriol

1323

326.6 proof of Lagrange’s theorem

We know that the cosets Hg form a partition of G (see the coset entry for proof of this.)

Since G is finite, we know it can be completely decomposed into a finite number of cosets.

Call this number n. We denote the ith coset by Hai and write G as

[ [ [

G = Ha1 Ha2 ··· Han

since each coset has |H| elements, we have

|G| = |H| · n

and so |H| divides |G|, which proves Lagrange’s theorem.

Version: 2 Owner: akrowne Author(s): akrowne

**326.7 proof of the converse of Lagrange’s theorem for
**

finite cyclic groups

**Following is a proof that, if G is a finite cyclic group and n ∈ Z+ is a divisor of |G|, then G
**

has a subgroup of order n.

**Let g be a generator of G. Then |g| = |hgi| = |G|. Let z ∈ Z such that nz = |G| = |g|.
**

|g|

Consider hg z i. Since g ∈ G, then g z ∈ G. Thus, hg z i ≤ G. Since |hg z i| = |g z | = GCD(z,|g|) =

nz nz z

GCD(z,nz)

= z = n, it follows that hg i is a subgroup of G of order n.

Version: 3 Owner: Wkbj79 Author(s): Wkbj79

326.8 proof that expG divides |G|

Following is a proof that exp G divides |G| for every finite group G.

**By the division algorithm, there exist q, r ∈ Z with 0 ≤ r < exp G such that |G| =
**

q(exp G) + r. Let g ∈ G. Then eG = g |G| = g q(exp G)+r = g q(exp G) g r = (g exp G )q g r =

(eG )q g r = eG g r = g r . Thus, for every g ∈ G, g r = eG . By the definition of exponent, r

cannot be positive. Thus, r = 0. It follows that exp G divides |G|.

Version: 4 Owner: Wkbj79 Author(s): Wkbj79

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326.9 proof that |g| divides expG

Following is a proof that, for every finite group G and for every g ∈ G, |g| divides exp G.

**By the division algorithm, there exist q, r ∈ Z with 0 ≤ r < |g| such that exp G = q|g| + r.
**

Since eG = g exp G = g q|g|+r = (g |g| )q g r = (eG )q g r = eG g r = g r , then, by definition of the

order of an element, r cannot be positive. Thus, r = 0. It follows that |g| divides exp G.

Version: 2 Owner: Wkbj79 Author(s): Wkbj79

326.10 proof that every group of prime order is cyclic

Following is a proof that every group of prime order is cyclic.

**Let p be a prime and G be a group such that |G| = p. Then G contains more than one
**

element. Let g ∈ G such that g 6= eG . Then hgi contains more than one element. Since

hgi ≤ G, then by Lagrange’s theorem, |hgi| divides p. Since |hgi| > 1 and |hgi| divides a

prime, then |hgi| = p = |G|. Hence, hgi = G. It follows that G is cyclic.

Version: 3 Owner: Wkbj79 Author(s): Wkbj79

1325

Chapter 327

20E05 – Free nonabelian groups

327.1 Nielsen-Schreier theorem

Let G be a free group and H a subgroup of G. Then H is free.

Version: 1 Owner: Evandar Author(s): Evandar

327.2 Scheier index formula

**Let G be a free group and H a subgroup of finite index |G : H| = n. By the Nielsen-Schreier theorem,
**

H is free. The Scheier index formula states that

rank(H) = n(rank(G) − 1) + 1.

**Thus implies more generally, if G0 is any group generated by m elements, then any subgroup
**

of index n can be generated by nm − n + 1 elements.

Version: 1 Owner: bwebste Author(s): bwebste

327.3 free group

**Let A be a set with elements ai for some index set I. We refer to A as an alphabet and the
**

elements of A as letters. A syllable is a symbol of the form ani for n ∈ Z. It is customary

to write a for a1 . Define a word to be a finite ordered string, or sequence, of syllables made

up of elements of A. For example,

a32 a1 a−1 2 −3

4 a3 a2

1326

is a five-syllable word. Notice that there exists a unique empty word, i.e. the word with

no syllables, usually written simply as 1. Denote the set of all words formed from elements

of A by W[A].

**Define a binary operation, called the product, on W[A] by concatenation of words. To
**

illustrate, if a32 a1 and a−1 4 3 −1 4

1 a3 are elements of W[A] then their product is simply a2 a1 a1 a3 .

This gives W[A] the structure of a semigroup with identity. The empty word 1 acts as a

right and left identity in W[A], and is the only element which has an inverse. In order to

give W[A] the structure of a group, two more ideas are needed.

**If v = u1 a0i u2 is a word where u1 , u2 are also words and ai is some element of A, an elemen-
**

tary contraction of type I replaces the occurrence of a0 by 1. Thus, after this type of

contraction we get another word w = u1 u2 . If v = u1 api aqi u2 is a word, an elementary con-

traction of type II replaces the occurrence of api aqi by aip+q which results in w = u1 ap+q

i u2 .

In either of these cases, we also say that w is obtained from v by an elementary contraction,

or that v is obtained from w by an elementary expansion.

**Call two words u, v equivalent (denoted u ∼ v) if one can be obtained from the other by
**

a finite sequence of elementary contractions or expansions. This is an equivalence relation

on W[A]. Let F[A] be the set of equivalence classes of words in W[A]. Then F[A] is group

under the operation

[u][v] = [uv]

where [u] ∈ F[A]. The inverse [u]−1 of an element [u] is obtained by reversing the order of

the syllables of [u] and changing the sign of each syllable. For example, if [u] = [a1 a23 ], then

[u]−1 = [a−2 −1

3 a1 ].

**We call F[A] the free group on the alphabet A or the free group generated by A.
**

A given group G is free if G is isomorphic to F[A] for some A. This seemingly ad hoc

construction gives an important result: Every group is the homomorphic image of some free

group.

Version: 4 Owner: jihemme Author(s): jihemme, rmilson, djao

**327.4 proof of Nielsen-Schreier theorem and Schreier
**

index formula

**While there are purely algebraic proofs of this fact, a much easier proof is available through
**

geometric group theory.

**Let G be a group which is free on a set X. Any group acts freely on its Cayley graph, and
**

the Cayley graph of G is a 2|X|-regular tree, which we will call T.

If H is any subgroup of G, then H also acts freely on T by restriction. Since groups that act freely on trees a

H is free.

1327

Moreover, we can obtain the rank of H (the size of the set on which it is free). If G is a finite

graph, then π1 (G) is free of rank −χ(G) − 1, where χ(G) denotes the Euler characteristic of

G. Since H ∼ = π1 (H\T), the rank of H is χ(H\T). If H is of finite index n in G, then H\T

is finite, and χ(H\T) = nχ(G\T). Of course −χ(G\T) − 1 is the rank of G. Substituting,

we find that

rank(H) = n(rank(G) − 1) + 1.

Version: 2 Owner: bwebste Author(s): bwebste

327.5 Jordan-Holder decomposition

A Jordan–Hölder decomposition of a group G is a filtration

G = G1 ⊃ G2 ⊃ · · · ⊃ Gn = {1}

**such that Gi+1 is a normal subgroup of Gi and the quotient Gi /Gi+1 is a simple group for
**

each i.

Version: 4 Owner: djao Author(s): djao

327.6 profinite group

**A topological group G is profinite if it is isomorphic to the inverse limit of some projective
**

system of finite groups. In other words, G is profinite if there exists a directed set I, a

collection of finite groups {Hi }i∈I , and homomorphisms αij : Hj → Hi for each pair i, j ∈ I

with i 6 j, satisfying

1. αii = 1 for all i ∈ I,

2. αij ◦ αjk = αik for all i, j, k ∈ I with i 6 j 6 k,

with the property that:

**• G is isomorphic as a group to the projective limit
**

( )

Y

lim Hi := (hi ) ∈ Hi αij (hj ) = hi for all i 6 j

←−

i∈I

under componentwise multiplication.

1328

Q

• The isomorphism from G to lim Hi (considered as a subspace of Hi ) is a homeomorphism

←− Q

of topological spaces, where each Hi is given the discrete topology and Hi is given

the product topology.

The topology on a profinite group is called the profinite topology.

Version: 3 Owner: djao Author(s): djao

327.7 extension

**A short exact sequence 0 → A → B → C → 0 is sometimes called an extension of C by A.
**

This term is also applied to an object B which fits into such an exact sequence.

Version: 1 Owner: bwebste Author(s): bwebste

327.8 holomorph

**Let K be a group, and let θ : Aut(K) → Aut(K) be the identity map. The holomorph of
**

K, denoted Hol(K), is the semidirect product K oθ Aut(K). Then K is a normal subgroup of

Hol(K), and any automorphism of K is the restriction of an inner automorphism of Hol(K).

For if φ ∈ Aut(K), then

(1, φ) · (k, 1) · (1, φ−1) = (1 · k θ(φ) , φ) · (1, φ−1 )

= (k θ(φ) · 1θ(φ) , φφ−1)

= (φ(k), 1).

Version: 2 Owner: dublisk Author(s): dublisk

**327.9 proof of the Jordan Holder decomposition theo-
**

rem

**Let |G| = N. We first prove existence, using induction on N. If N = 1 (or, more generally,
**

if G is simple) the result is clear. Now suppose G is not simple. Choose a maximal proper

normal subgroup G1 of G. Then G1 has a Jordan–Hölder decomposition by induction, which

produces a Jordan–Hölder decomposition for G.

1329

To prove uniqueness, we use induction on the length n of the decomposition series. If n = 1

then G is simple and we are done. For n > 1, suppose that

G ⊃ G1 ⊃ G2 ⊃ · · · ⊃ Gn = {1}

and

G ⊃ G01 ⊃ G02 ⊃ · · · ⊃ G0m = {1}

are two decompositions of G. If G1 = G01Tthen we’re done (apply the induction hypothesis

to G1 ), so assume G1 6= G01 . Set H := G1 G01 and choose a decomposition series

H ⊃ H1 ⊃ · · · ⊃ Hk = {1}

**for H. By the second isomorphism theorem, G1 /H = G1 G01 /G01 = G/G01 (the last equality
**

is because G1 G01 is a normal subgroup of G properly containing G1 ). In particular, H is a

normal subgroup of G1 with simple quotient. But then

G1 ⊃ G2 ⊃ · · · ⊃ Gn

and

G1 ⊃ H ⊃ · · · ⊃ H k

are two decomposition series for G1 , and hence have the same simple quotients by the

induction hypothesis; likewise for the G01 series. Therefore n = m. Moreover, since G/G1 =

G01 /H and G/G01 = G1 /H (by the second isomorphism theorem), we have now accounted for

all of the simple quotients, and shown that they are the same.

Version: 4 Owner: djao Author(s): djao

327.10 semidirect product of groups

**The goal of this exposition is to carefully explain the correspondence between the notions
**

of external and internal semi–direct products of groups, as well as the connection between

semi–direct products and short exact sequences.

Naturally, we start with the construction of semi–direct products.

**Definition 6. Let H and Q be groups and let θ : Q −→ Aut(H) be a group homomorphism.
**

The semi–direct product Hoθ Q is defined to be the group with underlying set {(h, q)such thath ∈

H, q ∈ Q} and group operation (h, q)(h0 , q 0 ) := (hθ(q)h0 , qq 0 ).

**We leave it to the reader to check that H oθ Q is really a group. It helps to know that the
**

inverse of (h, q) is (θ(q −1 )(h−1 ), q −1 ).

**For the remainder of this article, we omit θ from the notation whenever this map is clear
**

from the context.

1330

Set G := H o Q. There exist canonical monomorphisms H −→ G and Q −→ G, given by

h 7→ (h, 1Q ), h∈H

q 7→ (1H , q), q∈Q

**where 1H (resp. 1Q ) is the identity element of H (resp. Q). These monomorphisms are so
**

natural that we will treat H and Q as subgroups of G under these inclusions.

Theorem 3. Let G := H o Q as above. Then:

• H is a normal subgroup of G.

• HQ = G.

T

• H Q = {1G }.

**L et p : G −→ Q be the projection map defined by p(h, q) = q. Then p is a homomorphism
**

with kernel H. Therefore H is a normal subgroup of G.

Every (h, q) ∈ G can be written as (h, 1Q )(1H , q). Therefore HQ = G.

**Finally, it is evident that (1H , 1Q ) is the only element of G that is of the form (h, 1Q ) for
**

h ∈ H and (1H , q) for q ∈ Q.

This result motivates the definition of internal semi–direct products.

**Definition 7. Let G be a group with subgroups H and Q. We say G is the internal semi–
**

direct product of H and Q if:

• H is a normal subgroup of G.

• HQ = G.

T

• H Q = {1G }.

**We know an external semi–direct product is an internal semi–direct product (Theorem 3).
**

Now we prove a converse (Theorem 4), namely, that an internal semi–direct product is an

external semi–direct product.

T

Lemma 6. Let G be a group with subgroups H and Q. Suppose G = HQ and H Q = {1G }.

Then every element g of G can be written uniquely in the form hq, for h ∈ H and q ∈ Q.

**S ince G = HQ, we know that g can be T
**

written as hq. Suppose it can also be written as

0 −1

h q . Then hq = h q so h h = q q ∈ H Q = {1G }. Therefore h = h0 and q = q 0 .

0 0 0 0 0 −1

1331

Theorem 4. Suppose G is a group with subgroups H and Q, and G is the internal semi–

direct product of H and Q. Then G ∼

= H oθ Q where θ : Q −→ Aut(H) is given by

θ(q)(h) := qhq −1 , q ∈ Q, h ∈ H.

**B y lemma 6, every element g of G can be written uniquely in the form hq, with h ∈ H
**

and q ∈ Q. Therefore, the map φ : H o Q −→ G given by φ(h, q) = hq is a bijection from

G to H o Q. It only remains to show that this bijection is a homomorphism.

Given elements (h, q) and (h0 , q 0 ) in H o Q, we have

φ((h, q)(h0 , q 0 )) = φ((hθ(q)(h0 ), qq 0)) = φ(hqh0 q −1 , qq 0) = hqh0 q 0 = φ(h, q)φ(h0 , q 0 ).

Therefore φ is an isomorphism.

**Consider the external semi–direct product G := H oθ Q with subgroups H and Q. We know
**

from Theorem 4 that G is isomorphic to the external semi–direct product H oθ0 Q, where

we are temporarily writing θ0 for the conjugation map θ0 (q)(h) := qhq −1 of Theorem 4. But

in fact the two maps θ and θ0 are the same:

θ0 (q)(h) = (1H , q)(h, 1Q )(1H , q −1 ) = (θ(q)(h), 1Q ) = θ(q)(h).

**In summary, one may use Theorems 3 and 4 to pass freely between the notions of internal
**

semi–direct product and external semi–direct product.

**Finally, we discuss the correspondence between semi–direct products and split exact sequences
**

of groups.

Definition 8. An exact sequence of groups

i j

1 −→ H −→ G −→ Q −→ 1.

**is split if there exists a homomorphism k : Q −→ G such that j ◦ k is the identity map on Q.
**

Theorem 5. Let G, H, and Q be groups. Then G is isomorphic to a semi–direct product

H o Q if and only if there exists a split exact sequence

i j

1 −→ H −→ G −→ Q −→ 1.

F irst suppose G ∼

= H o Q. Let i : H −→ G be the inclusion map i(h) = (h, 1Q ) and let

j : G −→ Q be the projection map j(h, q) = q. Let the splitting map k : Q −→ G be the

inclusion map k(q) = (1H , q). Then the sequence above is clearly split exact.

**Now suppose we have the split exact sequence above. Let k : Q −→ G be the splitting map.
**

Then:

• i(H) = ker j, so i(H) is normal in G.

1332

• For any g ∈ G, set q := k(j(g)). Then j(gq −1 ) = j(g)j(k(j(g)))−1 = 1Q , so gq −1 ∈ Im i.

Set h := gq −1 . Then g = hq. Therefore G = i(H)k(Q).

**• Suppose g ∈ G is in both i(H) and k(Q). Write g = k(q). Then T k(q) ∈ Im i = ker j,
**

so q = j(k(q)) = 1Q . Therefore g = k(q) = k(1Q ) = 1G , so i(H) k(Q) = {1G }.

**This proves that G is the internal semi–direct product of i(H) and k(Q). These are iso-
**

morphic to H and Q, respectively. Therefore G is isomorphic to a semi–direct product

H o Q.

**Thus, not all normal subgroups H ⊂ G give rise to an (internal) semi–direct product G =
**

H o G/H. More specifically, if H is a normal subgroup of G, we have the canonical exact

sequence

1 −→ H −→ G −→ G/H −→ 1.

We see that G can be decomposed into H o G/H as an internal semi–direct product if and

only if the canonical exact sequence splits.

Version: 5 Owner: djao Author(s): djao

327.11 wreath product

Let A and B be groups, and let B act on the set Γ. Let AΓ be the set of all functions from

Γ to A. Endow AΓ with a group operation by pointwise multiplication. In other words, for

any f1 , f2 ∈ AΓ ,

(f1 f2 )(γ) = f1 (γ)f2 (γ) ∀γ ∈ Γ

where the operation on the right hand side above takes place in A, of course. Define the

action of B on AΓ by

bf (γ) := f (bγ),

for any f : Γ → A and all γ ∈ Γ.

**The wreath product of A and B according to the action of B on Γ, sometimes denoted
**

A oΓ B, is the following semidirect product of groups:

AΓ o B.

**Before going into further constructions, let us pause for a moment to unwind this definition.
**

Let W := A oΓ B. The elements of W are ordered pairs (f, b), for some function f : Γ → A

and some b ∈ B. The group operation in the semidirect product, for any (f1 , b1 ), (f2 , b2 ) ∈ W

is,

(f1 (γ), b1 )(f2 (γ), b2 ) = (f1 (γ)f2 (b1 γ), b1 b2 ) ∀γ ∈ Γ

The set AΓ can be interpreted as the cartesian product of A with itself, of cardinality Γ.

That is to say, Γ here plays the role of an index set for the Cartesian product. If Γ is finite,

1333

for instance, say Γ = {1, 2, . . . , n}, then any f ∈ AΓ is an n-tuple, and we can think of

any (f, b) ∈ W as the following ordered pair:

((a1 , a2 , . . . , an ), b) where a1 , a2 , . . . , an ∈ A

**The action of B on Γ in the semidirect product has the effect of permuting the elements of
**

the n-tuple f , and the group operation defined on AΓ gives pointwise multiplication. To be

explicit, suppose (f, a), (g, b) ∈ W , and for j ∈ Γ, f (j) = rj ∈ A and g(j) = sj ∈ A. Then,

(f, a)(g, b) = ((r1 , r2 , . . . , rn ), a)((s1 , s2 , . . . , sn ), b)

= ((r1 , r2 , . . . , rn )(sa1 , sa2 , . . . , san ), ab) (Notice the permutation of the indices!)

= ((r1 sa1 , r2 sa2 , . . . , rn san ), ab).

**A moment’s thought to understand this slightly messy notation will be illuminating (and
**

might also shed some light on the choice of terminology, “wreath” product).

Version: 11 Owner: bwebste Author(s): NeuRet

327.12 Jordan-Hlder decomposition theorem

Every finite group G has a filtration

G ⊃ G0 ⊃ · · · ⊃ Gn = {1},

**where each Gi+1 is normal in Gi and each quotient group Gi /Gi+1 is a simple group. Any two
**

such decompositions of G have the same multiset of simple groups Gi /Gi+1 up to ordering.

**A filtration of G satisfying the properties above is called a Jordan–Hölder decomposition of
**

G.

Version: 4 Owner: djao Author(s): djao

327.13 simplicity of the alternating groups

This is an elementary proof that for n > 5 the alternating group on n symbols, An , is simple.

**Throughout this discussion, fix n > 5. We will extensively employ cycle notation, with
**

composition on the left, as is usual. The following observation will also be useful. Let π be

a permutation written as disjoint cycles

π = (a1 , a2 , . . . , ak )(b1 , b2 , . . . , bl )(. . .) . . .

1334

It is easy to check that for any other permutation σ ∈ Sn

σπσ −1 = (σ(a1 ), σ(a2 ), . . . , σ(ak ))(σ(b1 ), σ(b2 ), . . .)(. . .) . . .)

**In particular, two permutations of Sn are conjugate exactly when they have the same cycle
**

type.

**Two preliminary results are necessary.
**

Lemma 7. An is generated by all cycles of length 3.

**A product of 3-cycles is an even permutation, so the subgroup generated by all 3-cycles is
**

therefore contained in An . For the reverse inclusion, by definition every even permutation

is the product of even number of transpositions. Thus, it suffices to show that the product

of two transpositions can be written as a product of 3-cycles. There are two possibilities.

Either the two transpositions move an element in common, say (a, b) and (a, c), or the two

transpositions are disjoint, say (a, b) and (c, d). In the former case,

(a, b)(a, c) = (a, c, b),

**and in the latter,
**

(a, b)(c, d) = (a, b, d)(c, b, d).

This establishes the first lemma.

Lemma 8. If a normal subgroup N / An contains a 3-cycle, then N = An .

**W e will show that if (a, b, c) ∈ N, then the assumption of normality implies that any other
**

(a0 , b0 , c0 ) ∈ N. This is easy to show, because there is some permutation in σ ∈ Sn that under

conjugation takes (a, b, c) to (a0 , b0 , c0 ), that is

σ(a, b, c)σ −1 = (σ(a), σ(b), σ(c)) = (a0 , b0 , c0 ).

**In case σ is odd, then (because n > 5) we can choose some transposition (d, e) ∈ An disjoint
**

from (a0 , b0 , c0 ) so that

σ(a, b, c)σ −1 = (d, e)(a0 , b0 , c0 )(d, e),

that is,

σ 0 (a, b, c)σ 0−1 = (d, e)σ(a, b, c)σ −1 (d, e) = (a0 , b0 , c0 )

where σ 0 is even. This means that N contains all 3-cycles, as N / An . Hence, by previous

lemma N = An as required.

**The rest of the proof proceeds by an exhuastive verification of all the possible cases. Suppose
**

there is some nontrivial N / An . We will show that N = An . In each case we will suppose N

contains a particular kind of element, and the normality will imply that N also contains a

certain conjugate of the element in An , thereby reducing the situation to a previously solved

case.

1335

Case 1 Suppose N contains a permutation π that when written as disjoint cycles has a

cycle of length at least 4, say

π = (a1 , a2 , a3 , a4 , . . .) . . .

Upon conjugation by (a1 , a2 , a3 ) ∈ An , we obtain

π 0 = (a1 , a2 , a3 )π(a3 , a2 , a1 ) = (a2 , a3 , a1 , a4 , . . .) . . .

**so that π 0 ∈ N, and also π 0 π −1 = (a1 , a2 , a4 ) ∈ N. Notice that the rest of the cycles cancel.
**

By Lemma 8, N = An .

**Case 2 The cyclic decompositions of elements of N only involve cycles of length 2 and at
**

least two cycles of length 3. Consider then π = (a, b, c)(d, e, f ) . . . Conjugation by (c, d, e)

implies that N also contains

π 0 = (c, d, e)π(e, d, c) = (a, b, d)(e, c, f ) . . . ,

and hence N also contains π 0 π = (a, d, c, b, f ) . . ., which reduces to Case 1.

**Case 3 There is an element of N whose cyclic decomposition only involves transpositions
**

and exactly one 3-cycle. Upon squaring, this element becomes a 3-cycle and Lemma 8

applies.

**Case 4 There is an element of N of the form π = (a, b)(c, d). Conjugating by (a, e, b) with
**

e distinct from a, b, c, d (again, at least one such e, as n > 5) yields

π 0 = (a, e, b)π(b, e, a) = (a, e)(c, d) ∈ N.

Hence π 0 π = (a, b, e) ∈ N. Lemma 8, applies and N = An .

**Case 5 Every element of N is the product of at least four transpositions. Suppose N
**

contains π = (a1 , b1 )(a2 , b2 )(a3 , b3 )(a4 , b4 ) . . ., the number of transpostions being even, of

course. This time we conjugate by (a2 , b1 )(a3 , b2 ).

π 0 = (a2 , b1 )(a3 , b2 )π(a3 , b2 )(a2 , b1 ) = (a1 , a2 )(a3 , b1 )(b2 , b3 )(a4 , b4 ),

and π 0 π = (a1 , a3 , b2 )(a2 , b3 , b1 ) ∈ N which is Case 2.

**Since this covers all possible cases, N = An and the alternating group contains no proper
**

nontrivial normal subgroups. QED.

Version: 8 Owner: rmilson Author(s): NeuRet

1336

327.14 abelian groups of order 120

Here we present an application of the fundamental theorem of finitely generated abelian groups.

Example (abelian groups of order 120):

**Let G be an abelian group of order n = 120. Since the group is finite it is obviously
**

finitely generated, so we can apply the theorem. There exist n1 , n2 , . . . , ns with

G∼

= Z/n1 Z ⊕ Z/n2 Z ⊕ . . . ⊕ Z/ns Z

∀i, ni > 2; ni+1 | ni for 1 6 i 6 s − 1

Notice that in the case of a finite group, r, as in the statement of the theorem, must be equal

to 0. We have s

Y

3

n = 120 = 2 · 3 · 5 = ni = n1 · n2 · . . . · ns

i=1

**and by the divisibility properties of ni we must have that every prime divisor of n must
**

divide n1 . Thus the possibilities for n1 are the following

2 · 3 · 5, 22 · 3 · 5, 23 · 3 · 5

**If n1 = 23 · 3 · 5 = 120 then s = 1. In the case that n1 = 22 · 3 · 5 then n2 = 2 and
**

s = 2. It remains to analyze the case n1 = 2 · 3 · 5. Now the possibilities for n2 are 2 (with

n3 = 2, s = 3) or 4 (with s = 2).

**Hence if G is an abelian group of order 120 it must be (up to isomorphism) one of the
**

following:

Z/120Z, Z/60Z ⊕ Z/2Z, Z/30Z ⊕ Z/4Z, Z/30Z ⊕ Z/2Z ⊕ Z/2Z

Also notice that they are all non-isomorphic. This is because

Z/(n · m)Z ∼

= Z/nZ ⊕ Z/mZ ⇔ gcd(n, m) = 1

which is due to the Chinese remainder theorem.

Version: 1 Owner: alozano Author(s): alozano

**327.15 fundamental theorem of finitely generated abelian
**

groups

**Theorem 2 (Fundamental Theorem of finitely generated abelian groups). Let G
**

be a finitely generated abelian group. Then there is a unique expression of the form:

G∼

= Zr ⊕ Z/n1 Z ⊕ Z/n2 Z ⊕ . . . ⊕ Z/ns Z

1337

for some integers r, ni satisfying:

r > 0; ∀i, ni > 2; ni+1 | ni for 1 6 i 6 s − 1

Version: 1 Owner: bwebste Author(s): alozano

327.16 conjugacy class

**Let G a group, and consider its operation (action) on itself give by conjugation, that is, the
**

mapping

(g, x) 7→ gxg −1

**Since conjugation is an equivalence relation, we obtain a partition of G into equivalence classes,
**

called conjugacy classes. So, the conjugacy class of X (represented Cx or C(x) is given by

Cx = {y ∈ X : y = gxg −1 for some g ∈ G}

Version: 2 Owner: drini Author(s): drini, apmxi

327.17 Frattini subgroup

**Let G be a group. The Frattini subgroup Φ(G) of G is the intersection of all maximal subgroups
**

of G.

Equivalently, Φ(G) is the subgroup of non-generators of G.

Version: 1 Owner: Evandar Author(s): Evandar

327.18 non-generator

**Let G be a group. g ∈ G is said to be a non-generator if whenever X is a generating set
**

for G then X r {g} is also a generating set for G.

Version: 1 Owner: Evandar Author(s): Evandar

1338

Chapter 328

**20Exx – Structure and classification
**

of infinite or finite groups

328.1 faithful group action

Let A be a G-set. That is, a set over which acts (or operates) a group G.

**The map mg : A → A defined as
**

mg (x) = ψ(g, x)

where g ∈ G and ψ is the action, is a permutation of A (in other words, a bijective function

of A) and so an element of SA . We can even get an homorphism from G to SA by the rule

g 7→ mg .

**If for any pair g, h ∈ G g 6= h we have mg 6= mh , in other words, the homomorphism g → mg
**

being injective, we say that the action is faithful.

Version: 3 Owner: drini Author(s): drini, apmxi

1339

Chapter 329

20F18 – Nilpotent groups

329.1 classification of finite nilpotent groups

Let G be a finite group. The following are equivalent:

1. G is nilpotent.

2. Every subgroup of G is subnormal.

3. Every subgroup H 6 G is properly contained in its normalizer.

4. Every maximal subgroup is normal.

5. Every Sylow subgroup is normal.

6. G is a direct product of p-groups.

Version: 1 Owner: bwebste Author(s): bwebste

329.2 nilpotent group

We define the lower central series of a group G to be the filtration of subgroups

G = G1 ⊃ G2 ⊃ · · ·

defined inductively by:

G1 := G,

Gi := [Gi−1 , G], i > 1,

1340

where [Gi−1 , G] denotes the subgroup of G generated by all commutators of the form hkh−1 k −1

where h ∈ Gi−1 and k ∈ G. The group G is said to be nilpotent if Gi = 1 for some i.

**Nilpotent groups can also be equivalently defined by means of upper central series. For a
**

group G, the upper central series of G is the filtration of subgroups

C1 ⊂ C2 ⊂ · · ·

**defined by setting C1 to be the center of G, and inductively taking Ci to be the unique
**

subgroup of G such that Ci /Ci−1 is the center of G/Ci−1 , for each i > 1. The group G is

nilpotent if and only if G = Ci for some i.

**Nilpotent groups are related to nilpotent Lie algebras in that a Lie group is nilpotent as
**

a group if and only if its corresponding Lie algebra is nilpotent. The analogy extends to

solvable groups as well: every nilpotent group is solvable, because the upper central series is

a filtration with abelian quotients.

Version: 3 Owner: djao Author(s): djao

1341

Chapter 330

**20F22 – Other classes of groups
**

defined by subgroup chains

330.1 inverse limit

Let {Gi }∞

i=0 be a sequence of groups which are related by a chain of surjective homomorphisms

fi : Gi → Gi−1 such that

f1 f2 f3 f4

G0 G1 G2 G3 ...

**Definition 28. The inverse limit of (Gi , fi ), denoted by
**

lim(Gi , fi ), or lim Gi

←− ←−

Q∞

is the subset of G

i=0 i formed by elements satisfying

( g0 , g1 , g2 , g3 , . . .), with gi ∈ Gi , fi (gi ) = gi−1

Q∞

Note: The inverse limit of Gi can be checked to be a subgroup of the product i=0 Gi . See

below for a more general definition.

Examples:

**1. Let p ∈ N be a prime. Let G0 = {0} and Gi = Z/pi Z. Define the connecting
**

homomorphisms fi , for i > 2, to be “reduction modulo pi−1 ” i.e.

fi : Z/pi Z → Z/pi−1 Z

fi (x mod pi ) = x mod pi−1

which are obviously surjective homomorphisms. The inverse limit of (Z/pi Z, fi ) is

called the p-adic integers and denoted by

Zp = lim Z/pi Z

←−

1342

2. Let E be an elliptic curve defined over C. Let p be a prime and for any natural number

n write E[n] for the n-torsion group, i.e.

E[n] = {Q ∈ E | n · Q = O}

In this case we define Gi = E[pi ], and

fi : E[pi ] → E[pi−1 ], fi (Q) = p · Q

The inverse limit of (E[pi ], fi ) is called the Tate module of E and denoted

Tp (E) = lim E[pi ]

←−

**The concept of inverse limit can be defined in far more generality. Let (S, 6) be a directed set
**

and let C be a category. Let {Gα }α∈S be a collection of objects in the category C and let

{fα,β : Gβ → Gα | α, β ∈ S, α 6 β}

be a collection of morphisms satisfying:

1. For all α ∈ S, fα,α = IdGα , the identity morphism.

**2. For all α, β, γ ∈ S such that α 6 β 6 γ, we have fα,γ = fα,β ◦ fβ,γ (composition of
**

morphisms).

Definition 29. The inverse limit of ({Gα }α∈S , {fα,β }), denoted by

lim(Gα , fα,β ), or lim Gα

←− ←−

Q

is defined to be the set of all (gα ) ∈ α∈S Gα such that for all α, β ∈ S

α 6 β ⇒ fα,β (gβ ) = gα

For a good example of this more general construction, see infinite Galois theory.

Version: 6 Owner: alozano Author(s): alozano

1343

Chapter 331

**20F28 – Automorphism groups of
**

groups

331.1 outer automorphism group

**The outer automorphism group of a group is the quotient of its automorphism group by
**

its inner automorphism group:

Out(G) = Aut(G)/Inn(G).

Version: 7 Owner: Thomas Heye Author(s): yark, apmxi

1344

Chapter 332

20F36 – Braid groups; Artin groups

332.1 braid group

**Consider two sets of n points on the complex plane C2 , of the form (1, 0), . . . , (n, 0), and
**

of the form (1, 1), . . . , (n, 1). We connect these two sets of points via a series of paths

fi : I → C2 , such that fi (t) 6= fj (t) for i 6= j and any t ∈ [0, 1]. Also, each fi may only

intersect the planes (0, z) and (1, z) for t = 0 and 1 respectively. Thus, the picture looks

like a bunch of strings connecting the two sets of points, but possibly tangled. The path

f = (f1 , . . . , fn ) determines a homotopy class f, where we require homotopies to satisfy the

same conditions on the fi . Such a homotopy class f is called a braid on n strands. We can

obtain a group structure on the set of braids on n strands as follows. Multiplication of two

strands f, g is done by simply following f first, then g, but doing each twice as fast. That is,

f · g is the homotopy class of the path

(

f (2t) if 0 6 t 6 1/2

fg =

g(2t − 1) if 1/2 6 t 6 1

**where f and g are representatives for f and g respectively. Inverses are done by following
**

the same strand backwards, and the identity element is the strand represented by straight

lines down. The result is known as the braid group on n strands, it is denoted by Bn .

**The braid group determines a homomorphism φ : Bn → Sn , where Sn is the symmetric group
**

on n letters. For f ∈ Bn , we get an element of Sn from map sending i 7→ p1 (fi (1)) where f

is a representative of the homtopy class f, and p1 is the projection onto the first factor. This

works because of our requirement on the points that the braids start and end, and since our

homotopies fix basepoints. The kernel of φ consists of the braids that bring each strand to

its original order. This kernel gives us the pure braid group on n strands, and is denoted

by Pn . Hence, we have a short exact sequence

1 → Pn → Bn → Sn → 1.

1345

We can also describe braid groups as certain fundamental groups, and in more generality.

Let M be a manifold, The configuration space of n ordered points on M is defined to

be Fn (M) = {(a1 , . . . , an ) ∈ M n | ai 6= aj fori 6= j}. The group Sn acts on Fn (M) by

permuting coordinates, and the corresponding quotient space Cn (M) = Fn (M)/Sn is called

the configuration space of n unordered points on M. In the case that M = C, we obtain the

regular and pure braid groups as π1 (Cn (M)) and π1 (Fn (M)) respectively.

**The group Bn can be given the following presentation. The presentation was given in Artin’s
**

first paper [1] on the braid group. Label the braids 1 through n as before. Let σi be the

braid that twists strands i and i + 1, with i passing beneath i + 1. Then the σi generate Bn ,

and the only relations needed are

σi σj = σj σi for |i − j| > 2, 1 6 i, j 6 n − 1

σi σi+1 σi = σi+1 σi σi+1 for 1 6 i 6 n − 2

**The pure braid group has a presentation with
**

−1 −1 −1

generatorsaij = σj−1 σj−2 · · · σi+1 σi2 σi+1 · · · σj−2 σj−1 for 1 6 i < j 6 n

**and defining relations
**

aij if i < r < s < j or r < s < i < j

a a a−1

rj ij rj if r<i=s<j

a−1

rs aij ars = −1 −1

arj asj aij asj arj if i=r<s<j

arj asj a−1 −1 −1 −1

rj asj aij asj arj asj arj if r<i<s<j

REFERENCES

1. E. Artin Theorie der Zöpfe. Abh. Math. Sem. Univ. Hamburg 4(1925), 42-72.

2. V.L. Hansen Braids and Coverings. London Mathematical Society Student Texts 18. Cambridge

University Press. 1989.

Version: 7 Owner: dublisk Author(s): dublisk

1346

Chapter 333

**20F55 – Reflection and Coxeter
**

groups

333.1 cycle

**Let S be a set. A cycle is a permutation (bijective function of a set onto itself) such that
**

there exist distinct elements a1 , a2 , . . . , ak of S such that

f (ai ) = ai+1 and f (ak ) = a1

that is

f (a1 ) = a2

f (a2 ) = a3

..

.

f (ak ) = a1

and f (x) = x for any other element of S.

**This can also be pictured as
**

a1 7→ a2 7→ a3 7→ · · · 7→ ak 7→ a1

and

x 7→ x

for any other element x ∈ S, where 7→ represents the action of f .

**One of the basic results on symmetric groups says that any permutation can be expressed
**

as product of disjoint cycles.

Version: 6 Owner: drini Author(s): drini

1347

333.2 dihedral group

**The nth dihedral group, Dn is the symmetry group of the regular n-sided polygon. The
**

group consists of n reflections, n − 1 rotations, and the identity transformation. Letting

ω = exp(2πi/n) denote a primitive nth root of unity, and assuming the polygon is centered

at the origin, the rotations Rk , k = 0, . . . , n − 1 (Note: R0 denotes the identity) are given

by

Rk : z 7→ ω k z, z ∈ C,

and the reflections Mk , k = 0, . . . , n − 1 by

Mk : z 7→ ω k z̄, z∈C

The abstract group structure is given by

Rk Rl = Rk+l , Rk Ml = Mk+l

Mk Ml = Rk−l , Mk Rl = Mk−l ,

where the addition and subtraction is carried out modulo n.

**The group can also be described in terms of generators and relations as
**

(M0 )2 = (M1 )2 = (M1 M0 )n = id.

This means that Dn is a rank-1 Coxeter group.

**Since the group acts by linear transformations
**

(x, y) → (x̂, ŷ), (x, y) ∈ R2

there is a corresponding action on polynomials p → p̂, defined by

p̂(x̂, ŷ) = p(x, y), p ∈ R[x, y].

The polynomials left invariant by all the group transformations form an algebra. This algebra

is freely generated by the following two basic invariants:

2 2 n n n−2 2

x +y , x − x y + ...,

2

the latter polynomial being the real part of (x + iy)n . It is easy to check that these two

polynomials are invariant. The first polynomial describes the distance of a point from the

origin, and this is unaltered by Euclidean reflections through the origin. The second polyno-

mial is unaltered by a rotation through 2π/n radians, and is also invariant with respect to

complex conjugation. These two transformations generate the nth dihedral group. Showing

that these two invariants polynomially generate the full algebra of invariants is somewhat

trickier, and is best done as an application of Chevalley’s theorem regarding the invariants

of a finite reflection group.

Version: 8 Owner: rmilson Author(s): rmilson

1348

Chapter 334

20F65 – Geometric group theory

334.1 groups that act freely on trees are free

**Let X be a tree, and Γ a group acting freely and faithfully by group automorphisms on X.
**

Then Γ is a free group.

**S ince Γ acts freely on X, the quotient graph X/Γ is well-defined, and X is the universal cover
**

of X/Γ since X is contractible. Thus Γ ∼ = π1 (X/Γ). Since any graph is homotopy equivalent

to a wedge of circles, and the fundamental group of such a space is free by Van Kampen’s theorem,

Γ is free.

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1349

Chapter 335

20F99 – Miscellaneous

335.1 perfect group

**A group G is called perfect if G = [G, G], where [G, G] is the derived subgroup of G, or
**

equivalently, if the abelianization of G is trivial.

Version: 1 Owner: bwebste Author(s): bwebste

1350

Chapter 336

**20G15 – Linear algebraic groups over
**

arbitrary fields

336.1 Nagao’s theorem

**For any integral domain k, the group of n×n invertible matrices with coefficients in k[t] is the
**

amalgamated free product of invertible matrices over k and invertible upper triangular matrices

over k[t], amalgamated over the upper triangular matrices of k. More compactly

GLn (k[t]) ∼

= GLn (k) ∗B(k) B(k[t]).

Version: 3 Owner: bwebste Author(s): bwebste

336.2 computation of the order of GL(n, Fq )

**GL(n, Fq ) is the group of n × n matrices over a finite field Fq with non-zero determinant.
**

Here is a proof that |GL(n, Fq )| = (q n − 1)(q n − q) · · · (q n − q n−1 ).

**Each element A ∈ GL(n, Fq ) is given by a collection of n Fq linearly independent vectors. If
**

one chooses the first column vector of A from (Fq )n there are q n choices, but one can’t choose

the zero vector since this would make the determinant of A zero. So there are really only

(q n − 1) choices. To choose an i-th vector from (Fq )n which is linearly independent from (i-1)

already choosen linearly independent vectors {V1 , · · · , Vi−1 } one must choose a vector not in

the span of {V1 , · · · , Vi−1 }. There are q i−1 vectors in this span so the number of choices is

clearly (q n − q i−1 ). Thus the number of linearly independent collections of n vectors in Fq

is: (q n − 1)(q n − q) · · · (q n − q n−1 ).

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336.3 general linear group

**Given a vector space V , the general linear group GL( V ) is defined to be the group of
**

invertible linear transformations from V to V . The group operation is defined by com-

position: given T : V −→ V and T 0 : V −→ V in GL( V ), the product T T 0 is just the

composition of the maps T and T 0 .

**If V = Fn for some field F, then the group GL(V ) is often denoted GL(n, F) or GLn (F).
**

In this case, if one identifies each linear transformation T : V −→ V with its matrix with

respect to the standard basis, the group GL(n, F) becomes the group of invertible n × n

matrices with entries in F, under the group operation of matrix multiplication.

Version: 3 Owner: djao Author(s): djao

**336.4 order of the general linear group over a finite
**

field

**GL(n, Fq ) is a finite group when Fq is a finite field with q elements. Furthermore, |GL(n, Fq )| =
**

(q n − 1)(q n − q) · · · (q n − q n−1 ).

Version: 16 Owner: benjaminfjones Author(s): benjaminfjones

336.5 special linear group

**Given a vector space V , the special linear group SL(V ) is defined to be the subgroup of the
**

general linear group GL(V ) consisting of all invertible linear transformations T : V −→ V

in GL(V ) that have determinant 1.

**If V = Fn for some field F, then the group SL(V ) is often denoted SL(n, F) or SLn (F), and if
**

one identifies each linear transformation with its matrix with respect to the standard basis,

then SL(n, F) consists of all n × n matrices with entries in F that have determinant 1.

Version: 2 Owner: djao Author(s): djao

1352

Chapter 337

**20G20 – Linear algebraic groups over
**

the reals, the complexes, the

quaternions

337.1 orthogonal group

**Let Q be a non-degenerate symmetric bilinear form over the real vector space Rn . A linear transformation
**

T : V −→ V is said to preserve Q if Q(T x, T y) = Q(x, y) for all vectors x, y ∈ V . The

subgroup of the general linear group GL(V ) consisting of all linear transformations that

preserve Q is called the orthogonal group with respect to Q, and denoted O(n, Q).

**If Q is also positive definite (i.e., Q is an inner product), then O(n, Q) is equivalent to the
**

group of invertible linear transformations that preserve the standard inner product on Rn ,

and in this case it is usually denoted O(n). One can show that a transformation T is in O(n)

if and only if T −1 = T T (the inverse of T equals the transpose of T ).

Version: 2 Owner: djao Author(s): djao

1353

Chapter 338

**20G25 – Linear algebraic groups over
**

local fields and their integers

338.1 Ihara’s theorem

**Let Γ be a discrete, torsion-free subgroup of SL2 Qp (where Qp is the field of p-adic numbers).
**

Then Γ is free.

**[ Proof, or a sketch thereof] There exists a p + 1 regular tree X on which SL2 Qp acts, with
**

stabilizer SL2 Zp (here, Zp denotes the ring T of p-adic integers). Since Zp is compact in its

profinite topology, so is SL2 Zp . Thus, SL2 Zp Γ must be compact, discrete and torsion-free.

Since compact and discrete implies finite, the only such group is trivial. Thus, Γ acts freely

on X. Since groups acting freely on trees are free, Γ is free.

Version: 6 Owner: bwebste Author(s): bwebste

1354

Chapter 339

**20G40 – Linear algebraic groups over
**

finite fields

339.1 SL2(F3)

**The special linear group over the finite field F3 is represented by SL2 (F3 ) and consists of the
**

2 × 2 invertible matrices with determinant equal to 1 and whose entries belong to F3 .

Version: 6 Owner: drini Author(s): drini, apmxi

1355

Chapter 340

20J06 – Cohomology of groups

340.1 group cohomology

**Let G be a group and let M be a (left) G-module. The 0th cohomology group of the
**

G-module M is

H 0 (G, M) = {m ∈ M : ∀σ ∈ G, σm = m}

which is the set of elements of M which are G-invariant, also denoted by M G .

**A map φ : G → M is said to be a crossed homomorphism (or 1-cocycle) if
**

φ(αβ) = φ(α) + αφ(β)

for all α, β ∈ G. If we fix m ∈ M, the map ρ : G → M defined by

ρ(α) = αm − m

is clearly a crossed homomorphism, said to be principal (or 1-coboundary). We define

the following groups:

Z 1 (G, M) = {φ : G → M : φ is a 1-cocycle}

B 1 (G, M) = {ρ : G → M : ρ is a 1-coboundary}

Finally, the 1st cohomology group of the G-module M is defined to be the quotient group:

H 1 (G, M) = Z 1 (G, M)/B 1 (G, M)

**The following proposition is very useful when trying to compute cohomology groups:
**

Proposition 1. Let G be a group and let A, B, C be G-modules related by an exact sequence:

0→A→B→C→0

Then there is a long exact sequence in cohomology:

0 → H 0 (G, A) → H 0 (G, B) → H 0 (G, C) → H 1 (G, A) → H 1 (G, B) → H 1(G, C)

1356

In general, the cohomology groups H n (G, M) can be defined as follows:

Definition 30. Define C 0 (G, M) = M and for n > 1 define the additive group:

C n (G, M) = {φ : Gn → M}

The elements of C n (G, M) are called n-cochains. Also, for n > 0 define the nth coboundary

homomorphism dn : C n (G, M) → C n+1 (G, M):

dn (f )(g1 , ..., gn+1 ) = g1 · f (g2, ..., gn+1 )

Xn

+ (−1)i f (g1 , ..., gi−1, gi gi+1 , gi+2 , ..., gn+1)

i=1

+ (−1)n+1 f (g1 , ..., gn )

Let Z n (G, M) = ker dn for n > 0, the set of n-cocyles. Also, let B 0 (G, M) = 1 and for n > 1

let B n (G, A) = image dn−1 , the set of n-coboundaries.

**Finally we define the nth -cohomology group of G with coefficients in M to be
**

H n (G, M) = Z n (G, M)/B n (G, M)

REFERENCES

1. J.P. Serre, Galois Cohomology, Springer-Verlag, New York.

2. James Milne, Elliptic Curves, online course notes.

3. Joseph H. Silverman, The Arithmetic of Elliptic Curves. Springer-Verlag, New York, 1986.

Version: 4 Owner: alozano Author(s): alozano

340.2 stronger Hilbert theorem 90

**Let K be a field and let K̄ be an algebraic closure of K. By K̄ + we denote the abelian group
**

(K̄, +) and similarly K̄ ∗ = (K̄, ∗) (here the operation is multiplication). Also we let

GK̄/K = Gal(K̄/K)

be the absolute Galois group of K.

Theorem 3 (Hilbert 90). Let K be a field.

1.

H 1 (GK̄/K , K̄ + ) = 0

1357

2.

H 1 (GK̄/K , K̄ ∗ ) = 0

3. If char(K), the characteristic of K, does not divide m (or char(K) = 0) then

H 1 (GK̄/K , µm ) ∼

= K ∗ /K ∗m

where µm denotes the set of all mth -roots of unity.

REFERENCES

1. J.P. Serre, Galois Cohomology, Springer-Verlag, New York.

2. J.P. Serre, Local Fields, Springer-Verlag, New York.

Version: 2 Owner: alozano Author(s): alozano

1358

Chapter 341

20J15 – Category of groups

341.1 variety of groups

**A variety of groups is the set of groups G such that all elements x1 , . . . , xn ∈ G satisfy a set
**

of equationally defined relations

ri (x1 , . . . , xn ) = 1 ∀i ∈ I,

where I is an index set.

**For example, abelian groups are a variety defined by the equations
**

{[x1 , x2 ] = 1},

where [x, y] = xyx−1 y −1.

**Nilpotent groups of class < c are a variety defined by
**

{[[· · · [[x1 , x2 ], x3 ] · · · ], xc ]}.

**Analogously, solvable groups of length < c are a variety. Abelian groups are a special case
**

of both of these.

Groups of exponent n are a variety, defined by {xn1 = 1}.

**A variety of groups is a full subcategory of the category of groups, and there is a free group
**

on any set of elements in the variety, which is the usual free group modulo the relations of the

variety applied to all elements. This satisfies the usual universal property of the free group

on groups in the variety, and is thus adjoint to the forgetful functor in the category of sets.

In the variety of abelian groups, we get back the usual free abelian groups. In the variety of

groups of exponent n, we get the Burnside groups.

Version: 1 Owner: bwebste Author(s): bwebste

1359

Chapter 342

20K01 – Finite abelian groups

342.1 Schinzel’s theorem

**Let a ∈ Q, not zero or 1 or −1. For any prime p which does not divide the numerator or
**

denominator of a in reduced form, a can be viewed as an element of the multiplicative group

Z/pZ. Let np be the order of this element in the multiplicative group.

Then the set of np over all such primes has finite complement in the set of positive integers.

One can generalize this as follows:

**Similarly, if K is a number field, choose a not zero or a root of unity in K. Then for any finite
**

place (discrete valuation) p with vp (a) = 0, we can view a as an element of the residue field

at p, and take the order np of this element in the multiplicative group.

Then the set of np over all such primes has finite complement in the set of positive integers.

Silverman also generalized this to elliptic curves over number fields.

References to come soon.

Version: 4 Owner: mathcam Author(s): Manoj, nerdy2

1360

Chapter 343

**20K10 – Torsion groups, primary
**

groups and generalized primary

groups

343.1 torsion

The

Definition 31. torsion of a group G is the set

Tor(G) = {g ∈ G : g n = e for some n ∈ N}.

A group is said to be

Definition 32. torsion-free if Tor(G) = {e}, i.e. the torsion consists only of the identity element.

If G is abelian then Tor(G) is a subgroup (the

**Definition 33. torsion group) of G.
**

Example 18 (Torsion of a cyclic group). For a cyclic group Zp , Tor(Zp ) = Zp .

In general, if G is a finite group then Tor(G) = G.

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1361

Chapter 344

**20K25 – Direct sums, direct products,
**

etc.

344.1 direct product of groups

**The external direct product G × H of two groups G and H is defined to be the set of
**

ordered pairs (g, h), with g ∈ G and h ∈ H. The group operation is defined by

(g, h)(g 0, h0 ) = (gg 0, hh0 )

**It can be shown that G × H obeys the group axioms. More generally, we can define the
**

external direct product of n groups, in the obvious way. Let G = G1 × . . . × Gn be the

set of all ordered n-tuples {(g1 , g2 . . . , gn ) | gi ∈ Gi } and define the group operation by

componentwise multiplication as before.

Version: 4 Owner: vitriol Author(s): vitriol

1362

Chapter 345

20K99 – Miscellaneous

345.1 Klein 4-group

**The Klein 4-group is the subgroup V (Vierergruppe) of S4 (see symmetric group) consisting
**

of the following 4 permutations:

(), (12), (34), (12)(34).

(see cycle notation). This is an abelian group, isomorphic to the product Z/2Z × Z/2Z. The

group is named after Felix Klein, a pioneering figure in the field of geometric group theory.

The Klein 4 group enjoys a number of interesting properties, some of which are listed below.

**1. It is the automorphism group of the graph consisting of two disjoint edges.
**

2. It is the unique 4 element group with the property that all elements are idempotent.

3. It is the symmetry group of a planar ellipse.

4. Consider the action of S4 , the permutation group of 4 elements, on the set of partitions

into two groups of two elements. There are 3 such partitions, which we denote by

(12, 34) (13, 24) (14, 23).

Thus, the action of S4 on these partition induces a homomorphism from S4 to S3 ; the

kernel is the Klein 4-group. This homomorphism is quite exceptional, and corresponds

to the fact that A4 (the alternating group) is not a simple group (notice that V is

actually a subgroup of A4 ). All other alternating groups are simple.

5. A more geometric way to see the above is the following: S4 is the group of symmetries

of a tetrahedron. There is an iduced action of S4 on the six edges of the tetrahedron.

Observing that this action preserves incidence relations one gets an action of S4 on the

three pairs of opposite edges (See figure).

1363

6. It is the symmetry group of the Riemannian curvature tensor.

3

4

1 2

Version: 7 Owner: rmilson Author(s): Dr Absentius, rmilson, imran

345.2 divisible group

**An abelian group D is said to be divisible if for any x ∈ D, n ∈ Z+ , there exists an element
**

x0 ∈ D such that nx0 = x.

Some noteworthy facts:

**• An abelian group is injective (as a Z-module) if and only if it is divisible.
**

• Every group is isomorphic to a subgroup of a divisible group.

**• Any divisible abelian group is isomorphic to the direct sum of its torsion subgroup and
**

n copies of the group of rationals (for some cardinal number n).

Version: 4 Owner: mathcam Author(s): mathcam

345.3 example of divisible group

**Let G denote the group of rational numbers taking the operation to be addition. Then for
**

any pq ∈ G and n ∈ Z+ , we have nq

p p

∈ G satisfying n nq = pq , so the group is divisible.

Version: 1 Owner: mathcam Author(s): mathcam

345.4 locally cyclic group

**A locally cyclic (or generalized cyclic) group is a group in which any pair of elements generates
**

a cyclic subgroup.

1364

Every locally cyclic group is abelian.

**If G is a locally cyclic group, then every finite subset of G generates a cyclic subgroup.
**

Therefore, the only finitely-generated locally cyclic groups are the cyclic groups themselves.

The group (Q, +) is an example of a locally cyclic group that is not cyclic.

Subgroups and quotients of locally cyclic groups are also locally cyclic.

A group is locally cyclic if and only if its lattice of subgroups is distributive.

Version: 10 Owner: yark Author(s): yark

1365

Chapter 346

20Kxx – Abelian groups

346.1 abelian group

**Let (G, ∗) be a group. If for any a, b ∈ G we have a ∗ b = b ∗ a, we say that the group is
**

abelian. Sometimes the expression commutative group is used, but this is less frequent.

**Abelian groups hold several interesting properties.
**

Theorem 4. If ϕ : G → G defined by ϕ(x) = x2 is a homomorphism, then G is abelian.

Proof. If such function were a homomorphism, we would have

(xy)2 = ϕ(xy) = ϕ(x)ϕ(y) = x2 y 2

**that is, xyxy = xxyy. Left-mutiplying by x−1 and right-multiplying by y −1 we are led to
**

yx = xy and thus the group is abelian. QED

Theorem 5. Any subgroup of an abelian group is normal.

**Proof. Let H be a subgroup of the abelian group G. Since ah = ha for any a ∈ G and any
**

h ∈ H we get aH = Ha. That is, H is normal in G. QED

Theorem 6. Quotient groups of abelian groups are also abelian.

**Proof Let H a subgroup of G. Since G is abelian, H is normal and we can get the quotient
**

group G/H whose elements are the equivalence classes for a ∼ b if ab−1 ∈ H.

**The operation on the quotient group is given by aH · bH = (ab)H. But bh · aH = (ba)H =
**

(ab)H, therefore the quotient group is also commutative. QED

Version: 12 Owner: drini Author(s): drini, yark, akrowne, apmxi

1366

Chapter 347

20M10 – General structure theory

347.1 existence of maximal semilattice decomposition

**Let S be a semigroup. A maximal semilattice decomposition for S is a surjective
**

homomorphism φ : S → Γ onto a semilattice Γ with the property that any other semilattice

decomposition factors through φ. So if φ0 : S → Γ0 is any other semilattice decomposition of

S, then there is a homomorphism Γ → Γ0 such that the following diagram commutes:

φ

S Γ

φ0

Γ0

Proposition 14. Every semigroup has a maximal semilattice decomposition.

**R ecall that each semilattice decompostion determines a semilattice congruence.
**

T If {ρi |

i ∈ I} is the family of all semilattice congruences on S, then define ρ = i∈I ρi . (Here, we

consider the congruences as subsets of S × S, and take their intersection as sets.)

**It is easy to see that ρ is also a semilattice congruence, which is contained in all other
**

semilattice congruences.

Therefore each of the homomorphisms S → S/ρi factors through S → S/ρ.

Version: 2 Owner: mclase Author(s): mclase

1367

347.2 semilattice decomposition of a semigroup

S

A semigroup S has a semilattice decomposition if we can write S = γ∈Γ Sγ as a disjoint

union of subsemigroups, indexed by elements of a semilattice Γ, with the additional condition

that x ∈ Sα and y ∈ Sβ implies xy ∈ Sαβ .

**Semilattice decompositions arise from homomorphims of semigroups onto semilattices. If
**

φ : S → Γ is a surjective homomorphism, then it is easy to see that we get a semilat-

tice decomposition by putting Sγ = φ−1 (γ) for each γ ∈ Γ. Conversely, every semilattice

decomposition defines a map from S to the indexing set Γ which is easily seen to be a

homomorphism.

**A third way to look at semilattice decompositions is to consider the congruence ρ defined
**

by the homomorphism φ : S → Γ. Because Γ is a semilattice, φ(x2 ) = φ(x) for all x, and

so ρ satisfies the constraint that x ρ x2 for all x ∈ S. Also, φ(xy) = φ(yx) so that xy ρ yx

for all x, y ∈ S. A congruence ρ which satisfies these two conditions is called a semilattice

congruence.

**Conversely, a semilattice congruence ρ on S gives rise to a homomorphism from S to a
**

semilattice S/ρ. The ρ-classes are the components of the decomposition.

Version: 3 Owner: mclase Author(s): mclase

347.3 simple semigroup

Let S be a semigroup. If S has no ideals other than itself, then S is said to be simple.

If S has no left ideals [resp. Right ideals] other than itself, then S is said to be left simple

[resp. right simple].

Right simple and left simple are stronger conditions than simple.

**A semigroup S is left simple if and only if Sa = S for all a ∈ S. A semigroup is both left
**

and right simple if and only if it is a group.

**If S has a zero element θ, then 0 = {θ} is always an ideal of S, so S is not simple (unless it
**

has only one element). So in studying semigroups with a zero, a slightly weaker definition is

required.

**Let S be a semigroup with a zero. Then S is zero simple, or 0-simple, if the following
**

conditions hold:

• S 2 6= 0

• S has no ideals except 0 and S itself

1368

The condition S 2 = 0 really only eliminates one semigroup: the 2-element null semigroup.

Excluding this semigroup makes parts of the structure theory of semigroups cleaner.

Version: 1 Owner: mclase Author(s): mclase

1369

Chapter 348

20M12 – Ideal theory

348.1 Rees factor

Let I be an ideal of a semigroup S. Define a congruence ∼ by x ∼ y iff x = y or x, y ∈ I.

**Then the Rees factor of S by I is the quotient S/ ∼. As a matter of notation, the
**

congruence ∼ is normally suppressed, and the quotient is simply written S/I.

**Note that a Rees factor always has a zero element. Intuitively, the quotient identifies all
**

element in I and the resulting element is a zero element.

Version: 1 Owner: mclase Author(s): mclase

348.2 ideal

**Let S be a semigroup. An ideal of S is a non-empty subset of S which is closed under
**

multiplication on either side by elements of S. Formally, I is an ideal of S if I is non-empty,

and for all x ∈ I and s ∈ S, we have sx ∈ I and xs ∈ I.

**One-sided ideals are defined similarly. A non-empty subset A of S is a left ideal (resp.
**

right ideal) of S if for all a ∈ A and s ∈ S, we have sa ∈ A (resp. as ∈ A).

**S a single element. If a 1∈ S, then the
**

A principal left ideal of S is a left ideal generated by

1

principal left ideal of S generated by a is S a = Sa {a}. (The notation S is explained

here.)

S

Similarly, the principal right ideal generated by a is aS 1 = aS {a}.

The notation L(a) and R(a) are also common for the principal left and right ideals generated

1370

by a respectively.

**A principal ideal of S is an ideal generated by a single element. The ideal generated by a
**

is [ [ [

S 1 aS 1 = SaS Sa aS {a}.

The notation J(a) = S 1 aS 1 is also common.

Version: 5 Owner: mclase Author(s): mclase

1371

Chapter 349

20M14 – Commutative semigroups

349.1 Archimedean semigroup

**Let S be a commutative semigroup. We say an element x divides an element y, written
**

x | y, if there is an element z such that xz = y.

**An Archimedean semigroup S is a commutative semigroup with the property that for all
**

x, y ∈ S there is a natural number n such that x | y n .

**This is related to the Archimedean property of positive real numbers R+ : if x, y > 0 then
**

there is a natural number n such that x < ny. Except that the notation is additive rather

than multiplicative, this is the same as saying that (R+ , +) is an Archimedean semigroup.

Version: 1 Owner: mclase Author(s): mclase

349.2 commutative semigroup

**A semigroup S is commutative if the defining binary operation is commutative. That is,
**

for all x, y ∈ S, the identity xy = yx holds.

**Although the term Abelian semigroup is sometimes used, it is more common simply to
**

refer to such semigroups as commutative semigroups.

A monoid which is also a commutative semigroup is called a commutative monoid.

Version: 1 Owner: mclase Author(s): mclase

1372

Chapter 350

20M20 – Semigroups of

transformations, etc.

350.1 semigroup of transformations

Let X be a set. A transformation of X is a function from X to X.

**If α and β are transformations on X, then their product αβ is defined (writing functions on
**

the right) by (x)(αβ) = ((x)α)β.

**With this definition, the set of all transformations on X becomes a semigroup, the full
**

semigroupf of transformations on X, denoted TX .

**More generally, a semigroup of transformations is any subsemigroup of a full set of
**

transformations.

**When X is finite, say X = {x1 , x2 , . . . , xn }, then the transformation α which maps xi to yi
**

(with yi ∈ X, of course) is often written:

x1 x2 . . . xn

α=

y1 y2 . . . yn

**With this notation it is quite easy to calculate products. For example, if X = {1, 2, 3, 4},
**

then

1 2 3 4 1 2 3 4 1 2 3 4

=

3 2 1 2 2 3 3 4 3 3 2 3

**When X is infinite, say X = {1, 2, 3, . . . }, then this notation is still useful for illustration
**

in cases where the transformation pattern is apparent. For example, if α ∈ TX is given by

1373

α : n 7→ n + 1, we can write

1 2 3 4 ...

α=

2 3 4 5 ...

Version: 3 Owner: mclase Author(s): mclase

1374

Chapter 351

20M30 – Representation of

semigroups; actions of semigroups on

sets

351.1 counting theorem

**Given a group action of a finite group G on a set X, the following expression gives the
**

number of distinct orbits

1 X

stabg (X)

|G| g∈G

Where stabg (X) is the number of elements fixed by the action of g.

Version: 8 Owner: mathcam Author(s): Larry Hammick, vitriol

351.2 example of group action

Let a, b, c be integers and let [a, b, c] denote the mapping

[a, b, c] : Z × Z 7→ Z, (x, y) 7→ ax2 + bxy + cy 2 .

**Let G be the group of 2 × 2 matrices such that det A = ±1 ∀ A ∈ G, and A ∈ G. The
**

substitution

txy 7→ A · txy

leads to

0 0 0

[a, b, c](a11 x + a12 y, a21x + a22 y) = a x2 + b xy + c y 2 ,

1375

where

0

a = a · a211 + b · a11 · a21 + c · a221 (351.2.1)

0

b = 2a · a11 · a12 + 2c · a21 · a22 + b(a11 a22 + a12 a21

0

c = a · a212 + b · a12 a22 + c · a222

So we define

0 0 0

[a, b, c] ∗ A := [a , b , c ]

0 0 0

to be the binary quadratic form with

coefficients a , b , c of x2 , xy, y 2, respectively as in

1 0

(495.2.1). Putting in A = we have [a, b, c] ∗ A = [a, b, c] for any binary quadratic

0 1

form [a, b, c]. Now let B be another matrix in G. We must show that

[a, b, c] ∗ (AB) = ([a, b, c] ∗ A) ∗ B.

00 00 00

Set [a, b, c] ∗ (AB) := [a , b , c ]. So we have

a · (a11 b11 + a12 b21 )2 + c · (a21 b11 + a22 b21 )2 + b · (a11 b11 + a12 b21 ) (a21 b11 + (351.2.2)

00

a = a22 b21 )

0 2 0 2

= a · b11 + c · b21 + (2a · a11 a12 + 2c · a21 a22 + b (a11 a22 + a12 a21 )) b11 b21

a · (a11 b12 + a12 b22 )2 + c · (a21 b12 + a22 b22 )2 + b · (a11 b12 + a12 b22 ) (a21 b12 + (351.2.3)

00

c = a22 b22 )

0 2 0 2

= a · b12 + c · b22 + (2a · a11 a12 + 2c · a21 a22 + b (a11 a22 + a12 a21 )) b12 b22

00

as desired. For the coefficient b we get

00

b = 2a · (a11 b11 + a12 b21 ) (a11 b12 + a12 b22 )

+ 2c · (a21 b11 + a22 b21 ) (a21 b12 + a22 b22 )

+ b · ((a11 b11 + a12 b21 ) (a21 b12 + a22 b22 ) + (a11 b12 + a12 b22 ) (a21 b11 + a22 b21 ))

and by evaluating the factors of b11 b12 , b21 b22 , and b11 b22 + b21 b12 , it can be checked that

00 0 0

b = 2a b11 b12 + 2c b21 b22 + (b11 b22 + b21 b12 ) (2a · a11 a12 + 2c · a21 a22 + b (a11 a22 + a12 a21 )) .

**This shows that
**

00 00 00 0 0 0

[a , b , c ] = [a , b , c ] ∗ B (351.2.4)

and therefore [a, b, c] ∗ (AB) = ([a, b, c] ∗ A) ∗ B. Thus, (495.2.1) defines an action of G on

the set of (integer) binary quadratic forms. Furthermore, the discriminant of each quadratic

form in the orbit of [a, b, c] under G is b2 − 4ac.

Version: 5 Owner: Thomas Heye Author(s): Thomas Heye

351.3 group action

**Let G be a group and let X be a set. A left group action is a function · : G × X −→ X such
**

that:

1376

1. 1G · x = x for all x ∈ X

2. (g1 g2 ) · x = g1 · (g2 · x) for all g1 , g2 ∈ G and x ∈ X

A right group action is a function · : X × G −→ X such that:

1. x · 1G = x for all x ∈ X

2. x · (g1 g2 ) = (x · g1 ) · g2 for all g1 , g2 ∈ G and x ∈ X

**There is a correspondence between left actions and right actions, given by associating the
**

right action x · g with the left action g · x := x · g −1. In many (but not all) contexts, it is

useful to identify right actions with their corresponding left actions, and speak only of left

actions.

Special types of group actions

**A left action is said to be effective, or faithful, if the function x 7→ g ·x is the identity function
**

on X only when g = 1G .

**A left action is said to be transitive if, for every x1 , x2 ∈ X, there exists a group element
**

g ∈ G such that g · x1 = x2 .

**A left action is free if, for every x ∈ X, the only element of G that stabilizes x is the identity;
**

that is, g · x = x implies g = 1G .

Faithful, transitive, and free right actions are defined similarly.

Version: 3 Owner: djao Author(s): djao

351.4 orbit

**Let G be a group, X a set, and · : G × X −→ X a group action. For any x ∈ X, the orbit
**

of x under the group action is the set

{g · x | g ∈ G} ⊂ X.

Version: 2 Owner: djao Author(s): djao

351.5 proof of counting theorem

**Let N be the cardinality of the set of all the couples (g, x) such that g · x = x. For each
**

g ∈ G, there exist stabg (X) couples with g as the first element, while for each x, there are

1377

|Gx | couples with x as the second element. Hence the following equality holds:

X X

N= stabg (X) = |Gx |.

g∈G x∈X

**From the orbit-stabilizer theorem it follows that:
**

X 1

N = |G| .

x∈X

|G(x)|

**Since all the x belonging to the same orbit G(x) contribute with
**

1

|G(x)| =1

|G(x)|

P

in the sum, then x∈X 1/|G(x)| precisely equals the number of distinct orbits s. We have

therefore X

stabg (X) = |G|s,

g∈G

which proves the theorem.

Version: 2 Owner: n3o Author(s): n3o

351.6 stabilizer

**Let G be a group, X a set, and · : G × X −→ X a group action. For any subset S of X, the
**

stabilizer of S, denoted Stab(S), is the subgroup

Stab(S) := {g ∈ G | g · s ∈ Sfor all s ∈ S}.

The stabilizer of a single point x in X is often denoted Gx .

Version: 3 Owner: djao Author(s): djao

1378

Chapter 352

20M99 – Miscellaneous

352.1 a semilattice is a commutative band

This note explains how a semilattice is the same as a commutative band.

**Let S be a semilattice, with partial order < and each pair of elements x and y hav-
**

ing a greatest lower bound x ∧ y. Then it is easy to see that the operation ∧ defines a

binary operation on S which makes it a commutative semigroup, and that every element is

idempotent since x ∧ x = x.

**Conversely, if S is such a semigroup, define x 6 y iff x = xy. Again, it is easy to see that
**

this defines a partial order on S, and that greatest lower bounds exist with respect to this

partial order, and that in fact x ∧ y = xy.

Version: 3 Owner: mclase Author(s): mclase

352.2 adjoining an identity to a semigroup

**It is possible to formally adjoin an identity element to any semigroup to make it into a
**

monoid.

S

Suppose S is a semigroup without an identity, and consider the set S S {1} where 1 is

a symbol not in S. Extend the semigroup operation from S to S {1} by additionally

defining: [

s · 1 = s = 1 · s, for alls ∈ S 1

It is easy to verify that this defines a semigroup (associativity is the only thing that needs

to be checked).

1379

As a matter of notation, it is customary to write S 1 for the semigroup S with an identity

adjoined in this manner, if S does not already have one, and to agree that S 1 = S, if S does

already have an identity.

**Despite the simplicity of this construction, however, it rarely allows one to simplify a problem
**

by considering monoids instead of semigroups. As soon as one starts to look at the structure

of the semigroup, it is almost invariably the case that one needs to consider subsemigroups

and ideals of the semigroup which do not contain the identity.

Version: 2 Owner: mclase Author(s): mclase

352.3 band

A band is a semigroup in which every element is idempotent.

A commutative band is called a semilattice.

Version: 1 Owner: mclase Author(s): mclase

352.4 bicyclic semigroup

**The bicyclic semigroup C(p, q) is the monoid generated by {p, q} with the single relation
**

pq = 1.

**The elements of C(p, q) are all words of the form q n pm for m, n > 0 (with the understanding
**

p0 = q 0 = 1). These words are multiplied as follows:

(

q n+k−m pl if m 6 k,

q n pm q k pl =

q n pl+m−k if m > k.

**It is apparent that C(p, q) is simple, for if q n pm is an element of C(p, q), then 1 = pn (q n pm )q m
**

and so S 1 q n pm S 1 = S.

It is useful to picture some further properties of C(p, q) by arranging the elements in a table:

1 p p2 p3 p4 ...

2 3

q qp qp qp qp4 ...

q 2 q 2 p q 2 p2 q 2 p3 q 2 p4 ...

q 3 q 3 p q 3 p2 q 3 p3 q 3 p4 ...

q 4 q 4 p q 4 p2 q 4 p3 q 4 p4 ...

.. .. .. .. .. ..

. . . . . .

1380

Then the elements below any horizontal line drawn through this table form a right ideal and

the elements to the right of any vertical line form a left ideal. Further, the elements on the

diagonal are all idempotents and their standard ordering is

1 > qp > q 2 p2 > q 3 p3 > · · · .

Version: 3 Owner: mclase Author(s): mclase

352.5 congruence

**Let S be a semigroup. An equivalence relation ∼ defined on S is called a congruence if it is
**

preserved under the semigroup operation. That is, for all x, y, z ∈ S, if x ∼ y then xz ∼ yz

and zx ∼ zy.

**If ∼ satisfies only x ∼ y implies xz ∼ yz (resp. zx ∼ zy) then ∼ is called a right congruence
**

(resp. left congruence).

Example 19. Suppose f : S → T is a semigroup homomorphism. Define ∼ by x ∼ y iff

f (x) = f (y). Then it is easy to see that ∼ is a congruence.

**If ∼ is a congruence, defined on a semigroup S, write [x] for the equivalence class of x
**

under ∼. Then it is easy to see that [x] · [y] = [xy] is a well-defined operation on the set of

equivalence classes, and that in fact this set becomes a semigroup with this operation. This

semigroup is called the quotient of S by ∼ and is written S/ ∼.

**Thus semigroup congruences are related to homomorphic images of semigroups in the same
**

way that normal subgroups are related to homomorphic images of groups. More precisely, in

the group case, the congruence is the coset relation, rather than the normal subgroup itself.

Version: 3 Owner: mclase Author(s): mclase

352.6 cyclic semigroup

A semigroup which is generated by a single element is called a cyclic semigroup.

Let S = hxi be a cyclic semigroup. Then as a set, S = {xn | n > 0}.

If all powers of x are distinct, then S = {x, x2 , x3 , . . . } is (countably) infinite.

**Otherwise, there is a least integer n > 0 such that xn = xm for some m < n. It is clear then
**

that the elements x, x2 , . . . , xn−1 are distinct, but that for any j ≥ n, we must have xj = xi

for some i, m ≤ i ≤ n − 1. So S has n − 1 elements.

1381

Unlike in the group case, however, there are in general multiple non-isomorphic cyclic semi-

groups with the same number of elements. In fact, there are t non-isomorphic cyclic semi-

groups with t elements: these correspond to the different choices of m in the above (with

n = t + 1).

The integer m is called the index of S, and n − m is called the period of S.

The elements K = {xm , xm+1 , . . . , xn−1 } are a subsemigroup of S. In fact, K is a cyclic group.

**A concrete representation of the semigroup with index m and period r as a semigroup of transformations
**

can be obtained as follows. Let X = {1, 2, 3, . . . , m + r}. Let

1 2 3 ... m + r −1 m + r

φ= .

2 3 4 ... m+r r+1

**Then φ generates a subsemigroup S of the full semigroup of transformations TX , and S is
**

cyclic with index m and period r.

Version: 3 Owner: mclase Author(s): mclase

352.7 idempotent

**An element x of a ring is called an idempotent element, or simply an idempotent if
**

x2 = x.

The set of idempotents of a ring can be partially ordered by putting e ≤ f iff e = ef = f e.

**The element 0 is a minimum element in this partial order. If the ring has an identity element,
**

1, then 1 is a maximum element in this partial order.

**Since these definitions refer only to the multiplicative structure of the ring, they also hold
**

for semigroups (with the proviso, of course, that a semigroup may not have a zero element).

In the special case of a semilattice, this partial order is the same as the one described in the

entry for semilattice.

**If a ring has an identity, then 1 − e is always an idempotent whenever e is an idempotent,
**

and e(1 − e) = (1 − e)e = 0.

**In a ring with an identity, two idempotents e and f are called a pair of orthogonal
**

idempotents if e + f = 1, and ef = f e = 0. Obviously, this is just a fancy way of saying

that f = 1 − e.

**More generally, a set {e1 , e2 , . . . , en } of idempotents is called a complete set of orthogonal
**

idempotents if ei ej = ej ei = 0 whenever i 6= j and if 1 = e1 + e2 + · · · + en .

1382

Version: 3 Owner: mclase Author(s): mclase

352.8 null semigroup

**A left zero semigroup is a semigroup in which every element is a left zero element. In
**

other words, it is a set S with a product defined as xy = x for all x, y ∈ S.

A right zero semigroup is defined similarly.

**Let S be a semigroup. Then S is a null semigroup if it has a zero element and if the
**

product of any two elements is zero. In other words, there is an element θ ∈ S such that

xy = θ for all x, y ∈ S.

Version: 1 Owner: mclase Author(s): mclase

352.9 semigroup

**A semigroup G is a set together with a binary operation · : G × G −→ G which satisfies the
**

associative property: (a · b) · c = a · (b · c) for all a, b, c ∈ G.

Version: 2 Owner: djao Author(s): djao

352.10 semilattice

**A lower semilattice is a partially ordered set S in which each pair of elements has a
**

greatest lower bound.

**A upper semilattice is a partially ordered set S in which each pair of elements has a least
**

upper bound.

**Note that it is not normally necessary to distinguish lower from upper semilattices, because
**

one may be converted to the other by reversing the partial order. It is normal practise to

refer to either structure as a semilattice and it should be clear from the context whether

greatest lower bounds or least upper bounds exist.

**Alternatively, a semilattice can be considered to be a commutative band, that is a semigroup
**

which is commutative, and in which every element is idempotent. In this context, semilattices

are important elements of semigroup theory and play a key role in the structure theory of

commutative semigroups.

A partially ordered set which is both a lower semilattice and an upper semilattice is a lattice.

1383

Version: 3 Owner: mclase Author(s): mclase

352.11 subsemigroup,, submonoid,, and subgroup

Let S be a semigroup, and let T be a subset of S.

**T is a subsemigroup of S if T is closed under the operation of S; that it if xy ∈ T for all
**

x, y ∈ T .

T is a submonoid of S if T is a subsemigroup, and T has an identity element.

T is a subgroup of S if T is a submonoid which is a group.

**Note that submonoids and subgroups do not have to have the same identity element as
**

S itself (indeed, S may not have an identity element). The identity element may be any

idempotent element of S.

**Let e ∈ S be an idempotent element. Then there is a maximal subsemigroup of S for which
**

e is the identity:

eSe = {exe | x ∈ S}.

In addition, there is a maximal subgroup for which e is the identity:

U(eSe) = {x ∈ eSe | ∃y ∈ eSe st xy = yx = e}.

**Subgroups with different identity elements are disjoint. To see this, suppose that G and H
**

are subgroups

T of a semigroup S with identity elements e and f respectively, and suppose

x ∈ G H. Then x has an inverse y ∈ G, and an inverse z ∈ H. We have:

e = xy = f xy = f e = zxe = zx = f.

Thus intersecting subgroups have the same identity element.

Version: 2 Owner: mclase Author(s): mclase

352.12 zero elements

**Let S be a semigroup. An element z is called a right zero [resp. left zero] if xz = z [resp.
**

zx = z] for all x ∈ S.

An element which is both a left and a right zero is called a zero element.

A semigroup may have many left zeros or right zeros, but if it has at least one of each, then

they are necessarily equal, giving a unique (two-sided) zero element.

1384

It is customary to use the symbol θ for the zero element of a semigroup.

Version: 1 Owner: mclase Author(s): mclase

1385

Chapter 353

**20N02 – Sets with a single binary
**

operation (groupoids)

353.1 groupoid

**A groupoid G is a set together with a binary operation · : G × G −→ G. The groupoid (or
**

“magma”) is closed under the operation.

There is also a separate, category-theoretic definition of “groupoid.”

Version: 7 Owner: akrowne Author(s): akrowne

353.2 idempotency

**If (S, ∗) is a magma, then an element x ∈ S is said to be idempotent if x ∗ x = x. If every
**

element of S is idempotent, then the binary operation ∗ (or the magma itself) is said to

be idempotent. For example, the ∧ and ∨ operations in a lattice are idempotent, because

x ∧ x = x and x ∨ x = x for all x in the lattice.

**A function f : D → D is said to be idempotent if f ◦ f = f . (This is just a special case
**

of the above definition, the magma in question being (D D , ◦), the monoid of all functions

from D to D, with the operation of function composition.) In other words, f is idempotent

iff repeated application of f has the same effect as a single application: f (f (x)) = f (x)

for all x ∈ D. An idempotent linear transformation from a vector space to itself is called a

projection.

Version: 12 Owner: yark Author(s): yark, Logan

1386

353.3 left identity and right identity

**Let G be a groupoid. An element e ∈ G is called a left identity element if ex = x for all
**

x ∈ G. Similarly, e is a right identity element if xe = x for all x ∈ G.

An element which is both a left and a right identity is an identity element.

**A groupoid may have more than one left identify element: in fact the operation defined by
**

xy = y for all x, y ∈ G defines a groupoid (in fact, a semigroup) on any set G, and every

element is a left identity.

**But as soon as a groupoid has both a left and a right identity, they are necessarily unique
**

and equal. For if e is a left identity and f is a right identity, then f = ef = e.

Version: 2 Owner: mclase Author(s): mclase

1387

Chapter 354

20N05 – Loops, quasigroups

354.1 Moufang loop

Proposition: Let Q be a nonempty quasigroup.

I) The following conditions are equivalent.

(x(yz))x = (xy)(zx) ∀x, y, z ∈ Q (354.1.1)

((xy)z)x = x(y(zx)) ∀x, y, z ∈ Q (354.1.2)

(yx)(zy) = (y(xz))y ∀x, y, z ∈ Q (354.1.3)

y(x(yz)) = ((yx)y)z ∀x, y, z ∈ Q (354.1.4)

II) If Q satisfies those conditions, then Q has an identity element (i.e. Q is a loop).

**For a proof, we refer the reader to the two references. Kunen in [1] shows that that any of
**

the four conditions implies the existence of an identity element. And Bol and Bruck [2] show

that the four conditions are equivalent for loops.

**Definition:A nonempty quasigroup satisfying the conditions (1)-(4) is called a Moufang
**

quasigroup or, equivalently, a Moufang loop (after Ruth Moufang, 1905-1977).

**The 16-element set of unit octonians over Z is an example of a nonassociative Moufang loop.
**

Other examples appear in projective geometry, coding theory, and elsewhere.

References

**[1] K. Kunen Moufang Quasigroups (PostScript format) (=Moufang Quasigroups, J. Al-
**

gebra 83 (1996) 231-234)

[2] R. H. Bruck, A Survey of Binary Systems, Springer-Verlag, 1958

1388

Version: 3 Owner: yark Author(s): Larry Hammick

354.2 loop and quasigroup

**A quasigroup is a groupoid G with the property that for every x, y ∈ G, there are unique
**

elements w, z ∈ G such that xw = y and zx = y.

A loop is a quasigroup which has an identity element.

**What distinguishes a loop from a group is that the former need not satisfy the associative
**

law.

Version: 1 Owner: mclase Author(s): mclase

1389

Chapter 355

**22-00 – General reference works
**

(handbooks, dictionaries,

bibliographies, etc.)

355.1 fixed-point subspace

**Let Σ ⊂ Γ be a subgroup where Γ is a compact Lie group acting on a vector space V . The
**

fixed-point subspace of Σ is

Fix(Σ) = {x ∈ V | σx = x, ∀σ ∈ Σ}

**Fix(Σ) is a linear subspace of V since
**

\

Fix(Σ) = ker(σ − I)

σ∈Σ

**where I is the identity. If it is important to specify the space V we use the following notation
**

FixV (Σ).

REFERENCES

[GSS] Golubsitsky, Matin. Stewart, Ian. Schaeffer, G. David: Singularities and Groups in Bifurca-

tion Theory (Volume II). Springer-Verlag, New York, 1988.

Version: 1 Owner: Daume Author(s): Daume

1390

Chapter 356

**22-XX – Topological groups, Lie
**

groups

356.1 Cantor space

**Cantor space denoted C is the set of all infinite binary sequences with the product topology.
**

It is a perfect Polish space. It is a compact subspace of Baire space, which is the set of all

infinite sequences of integers with the natural product topology.

REFERENCES

1. Moschovakis, Yiannis N. Descriptive set theory theory, 1980, Amsterdam ; New York : North-

Holland Pub. Co.

Version: 8 Owner: xiaoyanggu Author(s): xiaoyanggu

1391

Chapter 357

**22A05 – Structure of general
**

topological groups

357.1 topological group

**A topological group is a triple (G, ·, T) where (G, ·) is a group and T is a topology on
**

G such that under T, the group operation (x, y) 7→ x · y is continuous with respect to the

product topology on G × G and the inverse map x 7→ x−1 is continuous on G.

Version: 3 Owner: Evandar Author(s): Evandar

1392

Chapter 358

22C05 – Compact groups

358.1 n-torus

**The n-Torus, denoted T n , is a smooth orientable n dimensional manifold which is the product
**

1

of n 1-spheres, i.e. T n = S · · × S}1 .

| × ·{z

n

**Equivalently, the n-Torus can be considered to be Rn modulo the action (vector addition)
**

of the integer lattice Zn .

**The n-Torus is in addition a topological group. If we think of S 1 as the unit circle in C
**

1

and T n = S

| × ·{z· · × S}1 , then S 1 is a topological group and so is T n by coordinate-wise

n

multiplication. That is,

(z1 , z2 , . . . , zn ) · (w1 , w2 , . . . , wn ) = (z1 w1 , z2 w2 , . . . , zn wn )

Version: 2 Owner: ack Author(s): ack, apmxi

358.2 reductive

**Let G be a Lie group or algebraic group. G is called reductive over a field k if every
**

representation of G over k is completely reducible For example, a finite group is reduc-

tive over a field k if and only if its order is not divisible by the characteristic of k (by

Maschke’s theorem). A complex Lie group is reductive if and only if it is a direct product

of a semisimple group and an algebraic torus.

Version: 3 Owner: bwebste Author(s): bwebste

1393

Chapter 359

**22D05 – General properties and
**

structure of locally compact groups

359.1 Γ-simple

A representation V of Γ is Γ-simple if either

• V ∼

= W1 ⊕ W2 where W1 , W2 are absolutely irreducible for Γ and are Γ-isomorphic, or

• V is non-absolutely irreducible for Γ.

[GSS]

REFERENCES

[GSS] Golubsitsky, Matin. Stewart, Ian. Schaeffer, G. David.: Singularities and Groups in Bifur-

cation Theory (Volume II). Springer-Verlag, New York, 1988.

Version: 1 Owner: Daume Author(s): Daume

1394

Chapter 360

**22D15 – Group algebras of locally
**

compact groups

360.1 group C ∗-algebra

Let C[G] be the group ring of a discrete group G. It has two completions to a C ∗ -algebra:

**Reduced group C ∗ -algebra. The reduced group C ∗ -algebra, Cr∗ (G), is obtained by com-
**

pleting C[G] in the operator norm for its regular representation on l2 (G).

**Maximal group C ∗ -algebra. The maximal group C ∗ -algebra, Cmax ∗
**

(G) or just C ∗ (G), is

defined by the following universal property: any *-homomorphism from C[G] to some

B(H) (the C ∗ -algebra of bounded operators on some Hilbert space H) factors through

∗

the inclusion C[G] ,→ Cmax (G).

**If G is amenable then Cr∗ (G) ∼ ∗
**

= Cmax (G).

Version: 3 Owner: mhale Author(s): mhale

1395

Chapter 361

**22E10 – General properties and
**

structure of complex Lie groups

361.1 existence and uniqueness of compact real form

**Let G be a semisimple complex Lie group. Then there exists a unique (up to isomorphism)
**

real Lie group K such that K is compact and a real form of G. Conversely, if K is compact,

semisimple and real, it is the real form of a unique semisimple complex Lie group G. The

group K can be realized as the set of fixed points of a special involution of G, called the

Cartan involution.

**For example, the compact real form of SLn C, the complex special linear group, is SU(n), the
**

special unitary group. Note that SLn R is also a real form of SLn C, but is not compact.

**The compact real form of SOn C, the complex special orthogonal group, is SOn R, the real or-
**

thogonal group. SOn C also has other, non-compact real forms, called the pseudo-orthogonal

groups.

**The compact real form of Sp2n C, the complex symplectic group, is less well-known. It is
**

(unfortunately) also usually denoted Sp(2n), and consists of n × n “unitary” quaternion

matrices, that is,

Sp(2n) = {M ∈ GLn H|MM ∗ = I}

where M ∗ denotes M conjugate transpose. This different from the real symplectic group

Sp2n R.

Version: 2 Owner: bwebste Author(s): bwebste

1396

361.2 maximal torus

**Let K be a compact group, and let t ∈ K be an element whose centralizer has minimal
**

dimension (such elements are dense in K). Let T be the centralizer of t. This subgroup is

closed since T = ϕ−1 (t) where ϕ : K → K is the map k 7→ ktk −1 , and abelian since it is the

intersection of K with the Cartan subgroup of its complexification, and hence a torus, since

K (and thus T ) is compact. We call T a maximal torus of K.

**This term is also applied to the corresponding maximal abelian subgroup of a complex
**

semisimple group, which is an algebraic torus.

Version: 2 Owner: bwebste Author(s): bwebste

361.3 Lie group

**A Lie group is a group endowed with a compatible analytic structure. To be more precise,
**

Lie group structure consists of two kinds of data

**• a finite-dimensional, real-analytic manifold G
**

• and two analytic maps, one for multiplication G×G → G and one for inversion G → G,

which obey the appropriate group axioms.

**Thus, a homomorphism in the category of Lie groups is a group homomorphism that is
**

simultaneously an analytic mapping between two real-analytic manifolds.

**Next, we describe a natural construction that associates a certain Lie algebra g to every Lie
**

group G. Let e ∈ G denote the identity element of G.

**For g ∈ G let λg : G → G denote the diffeomorphisms corresponding to left multiplication
**

by g.

Definition 9. A vector-field V on G is called left-invariant if V is invariant with respect to

all left multiplications. To be more precise, V is left-invariant if and only if

(λg )∗ (V ) = V

(see push-forward of a vector-field) for all g ∈ G.

Proposition 15. The vector-field bracket of two left-invariant vector fields is again, a left-

invariant vector field.

**Proof. Let V1 , V2 be left-invariant vector fields, and let g ∈ G. The bracket operation is
**

covariant with respect to diffeomorphism, and in particular

(λg )∗ [V1 , V2 ] = [(λg )∗ V1 , (λg )∗ V2 ] = [V1 , V2 ].

1397

Q.E.D.

**Definition 10. The Lie algebra of G, denoted hereafter by g, is the vector space of all
**

left-invariant vector fields equipped with the vector-field bracket.

**Now a right multiplication is invariant with respect to all left multiplications, and it turns
**

out that we can characterize a left-invariant vector field as being an infinitesimal right mul-

tiplication.

**Proposition 16. Let a ∈ Te G and let V be a left-invariant vector-field such that Ve = a.
**

Then for all g ∈ G we have

Vg = (λg )∗ (a).

**The intuition here is that a gives an infinitesimal displacement from the identity element
**

and that Vg is gives a corresponding infinitesimal right displacement away from g. Indeed

consider a curve

γ : (−, ) → G

passing through the identity element with velocity a; i.e.

γ(0) = e, γ 0 (0) = a.

The above proposition is then saying that the curve

t 7→ gγ(t), t ∈ (−, )

passes through g at t = 0 with velocity Vg .

**Thus we see that a left-invariant vector-field is completely determined by the value it takes
**

at e, and that therefore g is isomorphic, as a vector space to Te G.

Of course, we can also consider the Lie algebra of right-invariant vector fields. The resulting

Lie-algebra is anti-isomorphic (the order in the bracket is reversed) to the Lie algebra of

left-invariant vector fields. Now it is a general principle that the group inverse operation

gives an anti-isomorphism between left and right group actions. So, as one may well expect,

the anti-isomorphism between the Lie algebras of left and right-invariant vector fields can

be realized by considering the linear action of the inverse operation on Te G.

**Finally, let us remark that one can induce the Lie algebra structure directly on Te G by
**

considering adjoint action of G on Te G.

Examples. [Coming soon.]

1398

Notes.

**1. No generality is lost in assuming that a Lie group has analytic, rather than C ∞ or
**

even C k , k = 1, 2, . . . structure. Indeed, given a C 1 differential manifold with a C 1

multiplication rule, one can show that the exponential mapping endows this manifold

with a compatible real-analytic structure.

Indeed, one can go even further and show that even C 0 suffices. In other words, a

topological group that is also a finite-dimensional topological manifold possesses a com-

patible analytic structure. This result was formulated by Hilbert as his fifth problem,

and proved in the 50’s by Montgomery and Zippin.

2. One can also speak of a complex Lie group, in which case G and the multiplication

mapping are both complex-analytic. The theory of complex Lie groups requires the

notion of a holomorphic vector-field. Not withstanding this complication, most of the

essential features of the real theory carry over to the complex case.

3. The name “Lie group” honours the Norwegian mathematician Sophus Lie who pio-

neered and developed the theory of continuous transformation groups and the corre-

sponding theory of Lie algebras of vector fields (the group’s infinitesimal generators,

as Lie termed them). Lie’s original impetus was the study of continuous symmetry of

geometric objects and differential equations.

The scope of the theory has grown enormously in the 100+ years of its existence. The

contributions of Elie Cartan and Claude Chevalley figure prominently in this evolution.

Cartan is responsible for the celebrated ADE classification of simple Lie algebras, as

well as for charting the essential role played by Lie groups in differential geometry and

mathematical physics. Chevalley made key foundational contributions to the analytic

theory, and did much to pioneer the related theory of algebraic groups. Armand Borel’s

book “Essays in the History of Lie groups and algebraic groups” is the definitive source

on the evolution of the Lie group concept. Sophus Lie’s contributions are the subject

of a number of excellent articles by T. Hawkins.

Version: 6 Owner: rmilson Author(s): rmilson

361.4 complexification

**Let G be a real Lie group. Then the complexification GC of G is the unique complex Lie
**

group equipped with a map ϕ : G → GC such that any map G → H where H is a complex Lie

group, extends to a holomorphic map GC → H. If g and gC are the respective Lie algebras,

gC ∼

= g ⊗R C.

**For simply connected groups, the construction is obvious: we simply take the simply con-
**

nected complex group with Lie algebra gC , and ϕ to be the map induced by the inclusion

g → gC .

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If γ ∈ G is central, then its image is in central in GC since g 7→ γgγ −1 is a map extending

ϕ, and thus must be the identity by uniqueness half of the universal property. Thus, if

Γ ⊂ G is a discrete central subgroup, then we get a map G/Γ → GC /ϕ(Γ), which gives a

complexification for G/Γ. Since every Lie group is of this form, this shows existence.

**Some easy examples: the complexification both of SLn R and SU(n) is SLn C. The complex-
**

ification of R is C and of S 1 is C∗ .

**The map ϕ : G → GC is not always injective. For example, if G is the universal cover
**

of SLn R (which has fundamental group Z), then GC ∼

= SLn C, and ϕ factors through the

covering G → SLn R.

Version: 3 Owner: bwebste Author(s): bwebste

361.5 Hilbert-Weyl theorem

theorem:

**Let Γ be a compact Lie group acting on V . Then there exists a finite Hilbert
**

basis for the ring P(Γ) (the set of invariant polynomials). [GSS]

proof:

In [GSS] on page 54.

theorem:(as stated by Hermann Weyl)

**The (absolute) invariants J(x, y, . . .) corresponding to a given set of representations
**

of a finite or a compact Lie group have a finite integrity basis. [PV]

proof:

In [PV] on page 274.

REFERENCES

[GSS] Golubsitsky, Matin. Stewart, Ian. Schaeffer, G. David.: Singularities and Groups in Bifur-

cation Theory (Volume II). Springer-Verlag, New York, 1988.

[HW] Hermann, Weyl: The Classical Groups: Their Invariants and Representations. Princeton

University Press, New Jersey, 1946.

Version: 3 Owner: Daume Author(s): Daume

1400

361.6 the connection between Lie groups and Lie alge-

bras

**Given a finite dimensional Lie group G, it has an associated Lie algebra g = Lie(G). The
**

Lie algebra encodes a great deal of information about the Lie group. I’ve collected a few

results on this topic:

**Theorem 7. (Existence) Let g be a finite dimensional Lie algebra over R or C. Then there
**

exists a finite dimensional real or complex Lie group G with Lie(G) = g.

**Theorem 8. (Uniqueness) There is a unique connected simply-connected Lie group G with
**

any given finite-dimensional Lie algebra. Every connected Lie group with this Lie algebra is

a quotient G/Γ by a discrete central subgroup Γ.

**Even more important, is the fact that the correspondence G 7→ g is functorial: given a
**

homomorphism ϕ : G → H of Lie groups, there is natural homomorphism defined on Lie

algebras ϕ∗ : g → h, which just the derivative of the map ϕ at the identity (since the Lie

algebra is canonically identified with the tangent space at the identity).

There are analogous existence and uniqueness theorems for maps:

**Theorem 9. (Existence) Let ψ : g → h be a homomorphism of Lie algebras. Then if G is
**

the unique connected, simply-connected group with Lie algebra g, and H is any Lie group

with Lie algebra h, there exists a homorphism of Lie groups ϕ : G → H with ϕ∗ = ψ.

**Theorem 10. (Uniqueness) Let G be connected Lie group and H an arbitrary Lie group.
**

Then if two maps ϕ, ϕ0 : G → H induce the same maps on Lie algebras, then they are equal.

**Essentially, what these theorems tell us is the correspondence g 7→ G from Lie algebras to
**

simply-connected Lie groups is functorial, and right adjoint to the functor H 7→ Lie(H) from

Lie groups to Lie algebras.

Version: 6 Owner: bwebste Author(s): bwebste

1401

Chapter 362

**26-00 – General reference works
**

(handbooks, dictionaries,

bibliographies, etc.)

362.1 derivative notation

This is the list of known standard representations and their nuances.

du df dy

, , −

dv dx dx

The most common notation, this is read as the derivative of u with respect

d2 y

to v. Exponents relate which derivative, for example, dx 2 is the second derivative of y with

resspect to x.

**f 0 (x) , f~0 (x) , y 00− This is read as f prime of x. The number of primes tells the derivative,
**

ie. f 000 (x) is the third derivative of f (x) with respect to x. Note that in higher dimensions,

this may be a tensor of a rank equal to the derivative.

**Dx f (x), Fy (x), fxy (x)− These notations are rather arcane, and should not be used gener-
**

ally, as they have other meanings. For example Fy can easily by the y component of a

vector-valued function. The subscript in this case means ”with respect to”, so Fyy would be

the second derivative of F with respect to y.

**D1 f (x), F2 (x), f12 (x)− The subscripts in these cases refer to the derivative with respect to
**

the nth variable. For example, F2 (x, y, z) would be the derivative of F with respect to y.

They can easily represent higher derivatives, ie. D21 f (x) is the derivative with respect to

the first variable of the derivative with respect to the second variable.

1402

∂u

∂v

, ∂f

∂x

− The partial derivative of u with respect to v. This symbol can be manipulated as

in du

dv

for higher partials.

d ∂

dv

, ∂v − This is the operator version of the derivative. Usually you will see it acting on

d

something such as dv (v 2 + 3u) = 2v.

**[Jf(x)] , [Df (x)]− The first of these represents the Jacobian of f, which is a matrix of partial
**

derivatives such that

D1 f1 (x) . . . Dn f1 (x)

.. .. ..

[Jf (x)] = . . .

D1 fm (x) . . . Dn fm (x)

where fn represents the nth function of a vector valued function. the second of these no-

tations represents the derivative matrix, which in most cases is the Jacobian, but in some

cases, does not exist, even though the Jacobian exists. Note that the directional derivative

in the direction ~v is simply [Jf(x)]~v .

Version: 7 Owner: slider142 Author(s): slider142

362.2 fundamental theorems of calculus

**The Fundamental Theorems of Calculus serve to demonstrate that integration and
**

differentiation are inverse processes.

First Fundamental Theorem:

**Suppose that F is a differentiable function on the interval [a, b]. Then intba F 0 (x) dx = F (b) −
**

F (a).

Second Fundamental Theorem:

**Let f be a continuous function on the interval [a, b], let c be an arbitrary point in this interval
**

and assume f is integrable on the intervals of the form [0, x] for all x ∈ [a, b]. Let F be defined

as F (x) = intxc f (t) dt for every x in (a, b). Then, F is differentiable and F 0 (x) = f (x).

**This result is about Riemann integrals. When dealing with Lebesgue integrals we get a
**

generalization with Lebesgue’s differentiation theorem.

Version: 9 Owner: mathcam Author(s): drini, greg

1403

362.3 logarithm

**Definition. Three real numbers x, y, p, with x, y > 0, are said to obey the logarithmic
**

relation

logx (y) = p

if they obey the corresponding exponential relation:

xp = y.

**Note that by the monotonicity and continuity property of the exponential operation, for
**

given x and y there exists a unique p satisfying the above relation. We are therefore able to

says that p is the logarithm of y relative to the base x.

Properties. There are a number of basic algebraic identities involving logarithms.

**logx (yz) = logx (y) + logx (z)
**

logx (y/z) = logx (y) − logx (z)

logx (y z ) = z logx (y)

logx (1) = 0

logx (x) = 1

logx (y) logy (x) = 1

logx (z)

logy (z) =

logx (y)

**Notes. In essence, logarithms convert multiplication to addition, and exponentiation to
**

multiplication. Historically, these properties of the logarithm made it a useful tool for doing

numerical calculations. Before the advent of electronic calculators and computers, tables of

logarithms and the logarithmic slide rule were essential computational aids.

**Scientific applications predominantly make use of logarithms whose base is the Eulerian number
**

e = 2.71828 . . .. Such logarithms are called natural logarithms and are commonly denoted

by the symbol ln, e.g.

ln(e) = 1.

Natural logarithms naturally give rise to the natural logarithm function.

**A frequent convention, seen in elementary mathematics texts and on calculators, is that
**

logarithms that do not give a base explicitly are assumed to be base 10, e.g.

log(100) = 2.

**This is far from universal. In Rudin’s “Real and Complex analysis”, for example, we see
**

a baseless log used to refer to the natural logarithm. By contrast, computer science and

1404

information theory texts often assume 2 as the default logarithm base. This is motivated

by the fact that log2 (N) is the approximate number of bits required to encode N different

messages.

The invention of logarithms is commonly credited to John Napier [ Biography]

Version: 13 Owner: rmilson Author(s): rmilson

**362.4 proof of the first fundamental theorem of calcu-
**

lus

**Let us make a subdivison of the intervalP[a, b], ∆ : {a = x0 < x1 < x2 < · · · < xn−1 < xn = b}
**

From this, we can say F (b) − F (a) = ni=1 [F (xi ) − F (xi−1 )].

**From the mean-value theorem, we have that for any two points, x̄ and x̄¯, ∃ ξ ∈ (x̄, x̄ ¯) 3
**

F (x̄¯) − F (x̄) = F (ξ)(x̄¯ − x̄) If we use xi as x̄¯ and xi−1 as x̄, calling our intermediate point

0

ξi , we get F (xi ) − F (xi−1 ) = F 0 (ξi )(xi − xi−1 ).

Combining

Pn these, and using the abbreviation ∆i x = xi − xi−1 , we have F (xi ) − F (xi−1 ) =

0

i=1 F (ξi )∆i xi .

P

From the definition of an integral ∀ > 0 ∃ δ > 0 3 | ni=1 F 0 (ξi )∆i x − intba F 0 (x) dx| <

when k∆k < δ. Thus, ∀ > 0, |F (b) − F (a) − intba F 0 (x) dx| < .

**lim→0 |F (b) − F (a) − intba F 0 (x) dx| = 0, but F (b) − F (a) − intba F 0 (x) dx is constant with
**

respect to , which can only mean that |F (b) − F (a) − int ba F 0 (x) dx| = 0, and so we have the

first fundamental theorem of calculus F (b) − F (a) = intba F 0 (x) dx.

Version: 4 Owner: greg Author(s): greg

**362.5 proof of the second fundamental theorem of cal-
**

culus

Recall that a continuous function is Riemann integrable, so the integral

F (x) = intxc f (t) dt

is well defined.

Consider the increment of F :

**F (x + h) − F (x) = intcx+h f (t) dt − intxc f (t) dt = intx+h
**

x f (t) dt

1405

(we have used the linearity of the integral with respect to the function and the additivity

with respect to the domain).

Now let M be the maximum of f on [x, x + h] and m be the minimum. Clearly we have

mh ≤ intxx+h f (t) dt ≤ Mh

**(this is due to the monotonicity of the integral with respect to the integrand) which can be
**

written as

F (x + h) − F (x) intxx+h f (t) dt

= ∈ [m, M]

h h

**Being f continuous, by the mean-value theorem, there exists ξh ∈ [x, x+h] such that f (ξh ) =
**

F (x+h)−F (x)

h

so that

F (x + h) − F (x)

F 0 (x) = lim = lim f (ξh ) = f (x)

h→0 h h→0

since ξh → x as h → 0.

Version: 1 Owner: paolini Author(s): paolini

362.6 root-mean-square

**If x1 , x2 , . . . , xn are real numbers, we define their root-mean-square or quadratic mean
**

as r

x21 + x22 + · · · + x2n

R(x1 , x2 , . . . , xn ) = .

n

**The root-mean-square of a random variable X is defined as the square root of the expectation
**

of X 2 : p

R(X) = E(X 2 )

**If X1 , X2 , . . . , Xn are random variables with standard deviations σ1 , σ2 , . . . , σn , then the stan-
**

dard deviation of their arithmetic mean, X1 +X2n+···+Xn , is the root-mean-square of σ1 , σ2 , . . . , σn .

Version: 1 Owner: pbruin Author(s): pbruin

362.7 square

The square of a number x is the number obtained multiplying x by itself. It’s denoted as x2 .

1406

Some examples:

52 = 25

2

1 1

=

3 9

02 = 0

.52 = .25

Version: 2 Owner: drini Author(s): drini

1407

Chapter 363

26-XX – Real functions

363.1 abelian function

**An abelian or hyperelliptic function is a generalisation of an elliptic function. It is a function
**

of two variables with four periods. In a similar way to an elliptic function it can also be

regarded as the inverse function to certain integrals (called abelian or hyperelliptic integrals)

of the form

dz

int

R(z)

where R is a polynomial of degree greater than 4.

Version: 2 Owner: vladm Author(s): vladm

363.2 full-width at half maximum

**The full-width at half maximum (FWHM) is a parameter used to describe the width of a
**

bump on a function (or curve). The FWHM is given by the distance beteen the points where

the function reaches half of its maximum value.

**For example: the function
**

10

f (x) = .

x2 +1

**f reaches its maximum for x = 0,(f (0) = 10), so f reaches half of its maximum value for
**

x = 1 and x = −1 (f (1) = f (−1) = 5). So the FWHM for f , in this case, is 2. Beacouse

the distance between A(1, 5) and B(−1, 5) si 2.

1408

The function

10

f (x) = .

x2

+1

is called ’The Agnesi curve’, from Maria Gaetana Agnesi (1718 - 1799).

Version: 2 Owner: vladm Author(s): vladm

1409

Chapter 364

**26A03 – Foundations: limits and
**

generalizations, elementary topology

of the line

364.1 Cauchy sequence

**A sequence x0 , x1 , x2 , . . . in a metric space (X, d) is a Cauchy sequence if, for every real number
**

> 0, there exists a natural number N such that d(xn , xm ) < whenever n, m > N.

Version: 4 Owner: djao Author(s): djao, rmilson

364.2 Dedekind cuts

**The purpose of Dedekind cuts is to provide a sound logical foundation for the real number
**

system. Dedekind’s motivation behind this project is to notice that a real number α, intu-

itively, is completely determined by the rationals strictly smaller than α and those strictly

larger than α. Concerning the completeness or continuity of the real line, Dedekind notes in

[2] that

If all points of the straight line fall into two classes such that every point of the

first class lies to the left of every point of the second class, then there exists one

and only one point which produces this division of all points into two classes,

this severing of the straight line into two portions.

**Dedekind defines a point to produce the division of the real line if this point is either the
**

least or greatest element of either one of the classes mentioned above. He further notes that

1410

the completeness property, as he just phrased it, is deficient in the rationals, which motivates

the definition of reals as cuts of rationals. Because all rationals greater than α are really just

excess baggage, we prefer to sway somewhat from Dedekind’s original definition. Instead,

we adopt the following definition.

Definition 34. A Dedekind cut is a subset α of the rational numbers Q that satisfies

these properties:

1. α is not empty.

2. Q \ α is not empty.

3. α contains no greatest element

4. For x, y ∈ Q, if x ∈ α and y < x, then y ∈ α as well.

**Dedekind cuts are particularly appealing for two reasons. First, they make it very easy to
**

prove the completeness, or continuity of the real line. Also, they make it quite plain to

distinguish the rationals from the irrationals on the real line, and put the latter on a firm

logical foundation. In the construction of the real numbers from Dedekind cuts, we make

the following definition:

Definition 35. A real number is a Dedekind cut. We denote the set of all real numbers

by R and we order them by set-theoretic inclusion, that is to say, for any α, β ∈ R,

α < β if and only if α ⊂ β

**where the inclusion is strict. We further define α = β as real numbers if α and β are equal
**

as sets. As usual, we write α 6 β if α < β or α = β. Moreover, a real number α is said to

be irrational if Q \ α contains no least element.

**The Dedekind completeness property of real numbers, expressed as the supremum property,
**

now becomes straightforward to prove. In what follows, we will reserve Greek variables for

real numbers, and Roman variables for rationals.

Theorem 11. Every nonempty subset of real numbers that is bounded above has a least upper bound.

**L et A be a nonempty set of real numbers, such that for every α ∈ A we have that α 6 γ
**

for some real number γ. Now define the set

[

sup A = α.

α∈A

We must show that this set is a real number. This amounts to checking the four conditions

of a Dedekind cut.

1411

1. sup A is clearly not empty, for it is the nonempty union of nonempty sets.

2. Because γ is a real number, there is some rational x that is not in γ. Since every α ∈ A

is a subset of γ, x is not in any α, so x 6∈ sup A either. Thus, Q \ sup A is nonempty.

3. If sup A had a greatest element g, then g ∈ α for some α ∈ A. Then g would be

a greatest element of α, but α is a real number, so by contrapositive, sup A has no

greatest element.

4. Lastly, if x ∈ sup A, then x ∈ α for some α, so given any y < x because α is a real

number y ∈ α, whence y ∈ sup A.

**Thus, sup A is a real number. Trivially, sup A is an upper bound of A, for every α ⊆ sup A.
**

It now suffices to prove that sup A 6 γ, because γ was an arbitrary upper bound. But this

is easy, because every x ∈ sup A is an element of α for some α ∈ A, so because α ⊆ γ, x ∈ γ.

Thus, sup A is the least upper bound of A. We call this real number the supremum of A.

**To finish the construction of the real numbers, we must endow them with algebraic opera-
**

tions, define the additive and multiplicative identity elements, prove that these definitions

give a field, and prove further results about the order of the reals (such as the totality of this

order) – in short, build a complete ordered field. This task is somewhat laborious, but we

include here the appropriate definitions. Verifying their correctness can be an instructive,

albeit tiresome, exercise. We use the same symbols for the operations on the reals as for the

rational numbers; this should cause no confusion in context.

Definition 36. Given two real numbers α and β, we define

**• The additive identity, denoted 0, is
**

0 := {x ∈ Q : x < 0}

**• The multiplicative identity, denoted 1, is
**

1 := {x ∈ Q : x < 1}

**• Addition of α and β denoted α + β is
**

α + β := {x + y : x ∈ α, y ∈ β}

**• The opposite of α, denoted −α, is
**

−α := {x ∈ Q : −x 6∈ α, but − x is not the least element of Q \ α}

**• The absolute value of α, denoted |α|, is
**

(

α, if α > 0

|α| :=

−α, if α 6 0

1412

• If α, β > 0, then multiplication of α and β, denoted α · β, is

α · β := {z ∈ Q : z 6 0 or z = xy for some x ∈ α, y ∈ β with x, y > 0}

In general,

0, if α = 0 or β = 0

α · β := |α| · |β| if α > 0, β > 0 or α < 0, β < 0

−(|α| · |β|) if α > 0, β < 0 or α > 0, β < 0

**• The inverse of α > 0, denoted α−1 , is
**

α−1 := {x ∈ Q : x 6 0 or x > 0 and (1/x) 6∈ α, but 1/x is not the least element of Q\α}

If α < 0,

α−1 := −(|α|)−1

**All that remains (!) is to check that the above definitions do indeed define a complete ordered
**

field, and that all the sets implied to be real numbers are indeed so. The properties of R

as an ordered field follow from these definitions and the properties of Q as an ordered field.

It is important to point out that in two steps, in showing that inverses and opposites are

properly defined, we require an extra property of Q, not merely in its capacity as an ordered

field. This requirement is the Archimedean property.

**Moreover, because R is a field of characteristic 0, it contains an isomorphic copy of Q. The
**

rationals correspond to the Dedekind cuts α for which Q \ α contains a least member.

REFERENCES

1. Courant, Richard and Robbins, Herbert. What is Mathematics? pp. 68-72 Oxford University

Press, Oxford, 1969

2. Dedekind, Richard. Essays on the Theory of Numbers Dover Publications Inc, New York 1963

3. Rudin, Walter Principles of Mathematical Analysis pp. 17-21 McGraw-Hill Inc, New York,

1976

4. Spivak, Michael. Calculus pp. 569-596 Publish or Perish, Inc. Houston, 1994

Version: 20 Owner: rmilson Author(s): rmilson, NeuRet

364.3 binomial proof of positive integer power rule

We will use the difference quotient in this proof of the power rule for positive integers. Let

f (x) = xn for some integer n > 0. Then we have

(x + h)n − xn

f 0 (x) = lim .

h→0 h

1413

We can use the binomial theorem to expand the numerator

C0n x0 hn + C1n x1 hn−1 + · · · + Cn−1

n

xn−1 h1 + Cnn xn h0 − xn

f 0 (x) = lim

h→0 h

n!

where Ckn = k!(n−k)!

. We can now simplify the above

0 hn + nxhn−1 + · · · + nxn−1 h + xn − xn

f (x) = lim

h→0 h

= lim (hn−1 + nxhn−2 + · · · + nxn−1 )

h→0

n−1

= nx

= nxn−1 .

Version: 4 Owner: mathcam Author(s): mathcam, slider142

364.4 exponential

**Preamble. We use R+ ⊂ R to denote the set of non-negative real numbers. Our aim is to
**

define the exponential, or the generalized power operation,

xp , x ∈ R+ , p ∈ R.

**The power p in the above expression is called the exponent. We take it as proven that R is
**

a complete, ordered field. No other properties of the real numbers are invoked.

**Definition. For x ∈ R+ and n ∈ Z we define xn in terms of repeated multiplication. To
**

be more precise, we inductively characterize natural number powers as follows:

x0 = 1, xn+1 = x · xn , n ∈ N.

**The existence of the reciprocal is guaranteed by the assumption that R is a field. Thus, for
**

negative exponents, we can define

x−n = (x−1 )n , n ∈ N,

where x−1 is the reciprocal of x.

**The case of arbitrary exponents is somewhat more complicated. A possible strategy is to
**

define roots, then rational powers, and then extend by continuity. Our approach is different.

For x ∈ R+ and p ∈ R, we define the set of all reals that one would want to be smaller than

xp , and then define the latter as the least upper bound of this set. To be more precise, let

x > 1 and define

L(x, p) = {z ∈ R+ : z n < xm for all m ∈ Z, n ∈ N such that m < pn}.

1414

We then define xp to be the least upper bound of L(x, p). For x < 1 we define

xp = (x−1 )p .

**The exponential operation possesses a number of important properties, some of which char-
**

acterize it up to uniqueness.

**Note. It is also possible to define the exponential operation in terms of the exponential function
**

and the natural logarithm. Since these concepts require the context of differential theory, it

seems preferable to give a basic definition that relies only on the foundational property of

the reals.

Version: 11 Owner: rmilson Author(s): rmilson

364.5 interleave sequence

**Let S be a set, and let {xi }, i = 0, 1, 2, . . . and {yi }, i = 0, 1, 2, . . . be two sequences in S.
**

The interleave sequence is defined to be the sequence x0 , y0 , x1 , y1 , . . . . Formally, it is the

sequence {zi }, i = 0, 1, 2, . . . given by

(

xk if i = 2k is even,

zi :=

yk if i = 2k + 1 is odd.

Version: 2 Owner: djao Author(s): djao

364.6 limit inferior

**Let S ⊂ R be a set of real numbers. Recall that a limit point of S is a real number x ∈ R
**

such that for all > 0 there exist infinitely many y ∈ S such that

|x − y| < .

**We define lim inf S, pronounced the limit inferior of S, to be the infimum of all the limit
**

points of S. If there are no limit points, we define the limit inferior to be +∞.

The two most common notations for the limit inferior are

lim inf S

and

lim S .

1415

An alternative, but equivalent, definition is available in the case of an infinite sequence of

real numbers x0 , x1 , x2 , , . . .. For each k ∈ N, let yk be the infimum of the k th tail,

yk = inf j>k xj .

This construction produces a non-decreasing sequence

y0 6 y1 6 y2 6 . . . ,

which either converges to its supremum, or diverges to +∞. We define the limit inferior of

the original sequence to be this limit;

lim inf xk = lim yk .

k k

Version: 7 Owner: rmilson Author(s): rmilson

364.7 limit superior

**Let S ⊂ R be a set of real numbers. Recall that a limit point of S is a real number x ∈ R
**

such that for all > 0 there exist infinitely many y ∈ S such that

|x − y| < .

We define lim sup S, pronounced the limit superior of S, to be the supremum of all the limit

points of S. If there are no limit points, we define the limit superior to be −∞.

**The two most common notations for the limit superior are
**

lim sup S

and

lim S .

**An alternative, but equivalent, definition is available in the case of an infinite sequence of
**

real numbers x0 , x1 , x2 , , . . .. For each k ∈ N, let yk be the supremum of the k th tail,

yk = sup xj .

j>k

**This construction produces a non-increasing sequence
**

y0 > y1 > y2 > . . . ,

which either converges to its infimum, or diverges to −∞. We define the limit superior of

the original sequence to be this limit;

lim sup xk = lim yk .

k k

Version: 7 Owner: rmilson Author(s): rmilson

1416

364.8 power rule

The power rule states that

D p

x = pxp−1 , p∈R

Dx

This rule, when combined with the chain rule, product rule, and sum rule, makes calculating

many derivatives far more tractable. This rule can be derived by repeated application of the

product rule. See the proof of the power rule.

Repeated use of the above formula gives

(

di k 0 i>k

x = k!

dxi (k−i)!

xk−i i 6 k,

for i, k ∈ Z.

Examples

D 0 0 D

x = =0= 1

Dx x Dx

D 1 D

x = 1x0 = 1 = x

Dx Dx

D 2

x = 2x

Dx

D 3

x = 3x2

Dx

D√ D 1/2 1 −1/2 1

x = x = x = √

Dx Dx 2 2 x

D e

2x = 2exe−1

Dx

Version: 4 Owner: mathcam Author(s): mathcam, Logan

364.9 properties of the exponential

The exponential operation possesses the following properties.

1417

• Homogeneity. For x, y ∈ R+ , p ∈ R we have

(xy)p = xp y p

• Exponent additivity. For x ∈ R+ we have

x0 = 1, x1 = x.

Furthermore

xp+q = xp xq , p, q ∈ R.

• Monotonicity. For x, y ∈ R+ with x < y and p ∈ R+ we have

xp < y p , x−p > y −p .

**• Continuity. The exponential operation is continuous with respect to its arguments.
**

To be more precise, the following function is continuous:

P : R+ × R → R, P (x, y) = xy .

**Let us also note that the exponential operation is characterized (in the sense of existence and
**

uniqueness) by the additivity and continuity properties. [Author’s note: One can probably

get away with substantially less, but I haven’t given this enough thought.]

Version: 10 Owner: rmilson Author(s): rmilson

364.10 squeeze rule

Squeeze rule for sequences

Let f, g, h : N → R be three sequences of real numbers such that

f (n) ≤ g(n) ≤ h(n)

**for all n. If limn→∞ f (n) and limn→∞ h(n) exist and are equal, say to a, then limn→∞ g(n)
**

also exists and equals a.

**The proof is fairly straightforward. Let e be any real number > 0. By hypothesis there exist
**

M, N ∈ N such that

|a − f (n)| < e for all n ≥ M

|a − h(n)| < e for all n ≥ N

Write L = max(M, N). For n ≥ L we have

1418

• if g(n) ≥ a:

|g(n) − a| = g(n) − a ≤ h(n) − a < e

• else g(n) < a and:

|g(n) − a| = a − g(n) ≤ a − f (n) < e

So, for all n ≥ L, we have |g(n) − a| < e, which is the desired conclusion.

Squeeze rule for functions

**Let f, g, h : S → R be three real-valued functions on a neighbourhood S of a real number b,
**

such that

f (x) ≤ g(x) ≤ h(x)

for all x ∈ S − {b}. If limx→b f (x) and limx→b h(x) exist and are equal, say to a, then

limx→b g(x) also exists and equals a.

Again let e be an arbitrary positive real number. Find positive reals α and β such that

|a − f (x)| < e whenever 0 < |b − x| < α

|a − h(x)| < e whenever 0 < |b − x| < β

Write δ = min(α, β). Now, for any x such that |b − x| < δ, we have

• if g(x) ≥ a:

|g(x) − a| = g(x) − a ≤ h(x) − a < e

• else g(x) < a and:

|g(x) − a| = a − g(x) ≤ a − f (x) < e

and we are done.

Version: 1 Owner: Daume Author(s): Larry Hammick

1419

Chapter 365

26A06 – One-variable calculus

365.1 Darboux’s theorem (analysis)

**Let f : [a, b] → R be a real-valued continuous function on [a, b], which is differentiable on
**

(a, b), differentiable from the right at a, and differentiable from the left at b. Then f 0 satisfies

the intermediate value theorem: for every t between f+0 (a) and f−0 (b), there is some x ∈ [a, b]

such that f 0 (x) = t.

**Note that when f is continuously differentiable (f ∈ C 1 ([a, b])), this is trivially true by the
**

intermediate value theorem. But even when f 0 is not continuous, Darboux’s theorem places

a severe restriction on what it can be.

Version: 3 Owner: mathwizard Author(s): mathwizard, ariels

365.2 Fermat’s Theorem (stationary points)

**Let f : (a, b) → R be a continuous function and suppose that x0 ∈ (a, b) is a local extremum
**

of f . If f is differentiable in x0 then f 0 (x0 ) = 0.

Version: 2 Owner: paolini Author(s): paolini

1420

365.3 Heaviside step function

**The Heaviside step function is the function H : R → R defined as
**

0 when x < 0,

H(x) = 1/2 when x = 0,

1 when x > 0.

Here, there are many conventions for the value at x = 0. The motivation for setting H(0) =

1/2 is that we can then write H as a function of the signum function (see this page). In

applications, such as the Laplace transform, where the Heaviside function is used extensively,

the value of H(0) is irrelevant.

**The function is named after Oliver Heaviside (1850-1925) [1]. However, the function was
**

already used by Cauchy[2], who defined the function as

1 √

u(t) = t + t/ t2

2

and called it a coefficient limitateur [1].

REFERENCES

1. The MacTutor History of Mathematics archive, Oliver Heaviside.

2. The MacTutor History of Mathematics archive, Augustin Louis Cauchy.

3. R.F. Hoskins, Generalised functions, Ellis Horwood Series: Mathematics and its applica-

tions, John Wiley & Sons, 1979.

Version: 1 Owner: Koro Author(s): matte

365.4 Leibniz’ rule

**Theorem [Leibniz’ rule] ([1] page 592) Let f and g be real (or complex) valued functions
**

that are defined on an open interval of R. If f and g are k times differentiable, then

Xk

(k) k (k−r) (r)

(f g) = f g .

r=0

r

For multi-indices, Leibniz’ rule have the following generalization:

**Theorem [2] If f, g : Rn → C are smooth functions, and j is a multi-index, then
**

X j

j

∂ (f g) = ∂ i (f ) ∂ j−i (g),

i≤j

i

where i is a multi-index.

1421

REFERENCES

1. R. Adams, Calculus, a complete course, Addison-Wesley Publishers Ltd, 3rd ed.

2. http://www.math.umn.edu/ jodeit/course/TmprDist1.pdf

Version: 3 Owner: matte Author(s): matte

365.5 Rolle’s theorem

**Rolle’s theorem. If f is a continuous function on [a, b], such that f (a) = f (b) = 0 and
**

differentiable on (a, b) then there exists a point c ∈ (a, b) such that f 0 (c) = 0.

Version: 8 Owner: drini Author(s): drini

365.6 binomial formula

The binomial formula gives the power series expansion of the pth power function for every

real power p. To wit,

X∞ n

p n x

(1 + x) = p , x ∈ R, |x| < 1,

n=0

n!

where

pn = p(p − 1) . . . (p − n + 1)

denotes the nth falling factorial of p.

**Note that for p ∈ N the power series reduces to a polynomial. The above formula is therefore
**

a generalization of the binomial theorem.

Version: 4 Owner: rmilson Author(s): rmilson

365.7 chain rule

**Let f (x), g(x) be differentiable, real-valued functions. The derivative of the composition
**

(f ◦ g)(x) can be found using the chain rule, which asserts that:

(f ◦ g)0(x) = f 0 (g(x)) g 0(x)

**The chain rule has a particularly suggestive appearance in terms of the Leibniz formalism.
**

Suppose that z depends differentiably on y, and that y in turn depends differentiably on x.

1422

Then,

dz dz dy

=

dx dy dx

The apparent cancellation of the dy term is at best a formal mnemonic, and does not con-

stitute a rigorous proof of this result. Rather, the Leibniz format is well suited to the

interpretation of the chain rule in terms of related rates. To wit:

**The instantaneous rate of change of z relative to x is equal to the rate of change
**

of z relative to y times the rate of change of y relative to x.

Version: 5 Owner: rmilson Author(s): rmilson

365.8 complex Rolle’s theorem

**Theorem [1] Suppose Ω is an open convex set in C, suppose f is a holomorphic function
**

f : Ω → C, and suppose f (a) = f (b) = 0 for distinct points a, b in Ω. Then there exist

points u, v on Lab (the straight line connecting a and b not containing the endpoints), such

that

Re{f 0 (u)} = 0 and Im{f 0(v)} = 0.

REFERENCES

1. J.-Cl. Evard, F. Jafari, A Complex Rolle’s Theorem, American Mathematical Monthly,

Vol. 99, Issue 9, (Nov. 1992), pp. 858-861.

Version: 4 Owner: matte Author(s): matte

365.9 complex mean-value theorem

**Theorem [1] Suppose Ω is an open convex set in C, suppose f is a holomorphic function
**

f : Ω → C, and suppose a, b are distinct points in Ω. Then there exist points u, v on Lab

(the straight line connecting a and b not containing the endpoints), such that

f (b) − f (a)

Re{ } = Re{f 0 (u)},

b−a

f (b) − f (a)

Im{ } = Im{f 0 (v)},

b−a

1423

REFERENCES

1. J.-Cl. Evard, F. Jafari, A Complex Rolle’s Theorem, American Mathematical Monthly,

Vol. 99, Issue 9, (Nov. 1992), pp. 858-861.

Version: 2 Owner: matte Author(s): matte

365.10 definite integral

**The definite integral with respect to x of some function f (x) over the closed interval [a, b]
**

is defined to be the “area under the graph of f (x) with respect to x” (if f(x) is negative,

then you have a negative area). It is written as:

intba f (x) dx

one way to find the value of the integral is to take a limit of an approximation technique as

the precision increases to infinity.

**For example, use a Riemann sum which approximates the area by dividing it into n intervals
**

of equal widths, and then calculating the area of rectangles with the width of the interval

and height dependent on the function’s value in the interval. Let Rn be this approximation,

which can be written as n

X

Rn = f (x∗i )∆x

i=1

where x∗i is some x inside the ith interval.

**Then, the integral would be
**

n

X

intba f (x) dx = lim Rn = lim f (x∗i )∆x

n→∞ n→∞

i=1

**We can use this definition to arrive at some important properties of definite integrals (a, b,
**

c are constant with respect to x):

intba f (x) + g(x) dx = intba f (x) dx + intba g(x) dx

intba f (x) − g(x) dx = intba f (x) dx − intba g(x) dx

intba f (x) dx = −intab f (x) dx

intba f (x) dx = intca f (x) dx + intbc f (x) dx

intba cf (x) dx = cintba f (x) dx

There are other generalisations about integrals, but many require the fundamental theorem of calculus.

Version: 4 Owner: xriso Author(s): xriso

1424

365.11 derivative of even/odd function (proof )

**Suppose f (x) = ±f (−x). We need to show that f 0 (x) = ∓f 0 (−x). To do this, let us
**

define the auxiliary function m : R → R, m(x) = −x. The condition on f is then f (x) =

±(f ◦ m)(x). Using the chain rule, we have that

f 0 (x) = ±(f ◦ m)0 (x)

= ±f 0 m(x) m0 (x)

= ∓f 0 (−x),

and the claim follows. 2

Version: 2 Owner: mathcam Author(s): matte

365.12 direct sum of even/odd functions (example)

Example. direct sum of even and odd functions

Let us define the sets

**F = {f |f is a function fromR toR},
**

F+ = {f ∈ F |f (x) = f (−x)for allx ∈ R},

F− = {f ∈ F |f (x) = −f (−x)for allx ∈ R}.

**In other words, F contain all functions from R to R, F+ ⊂ F contain all even functions, and
**

F− ⊂ F contain all odd functions. All of these spaces have a natural vector space structure:

for functions f and g we define f + g as the function x 7→ f (x) + g(x). Similarly, if c is a

real constant, then cf is the function x 7→ cf (x). With these operations, the zero vector is

the mapping x 7→ 0.

We claim that F is the direct sum of F+ and F− , i.e., that

F = F+ ⊕ F− . (365.12.1)

**To prove this claim, let us first note that F± are vector subspaces of F . Second, given an
**

arbitrary function f in F , we can define

1

f+ (x) = f (x) + f (−x) ,

2

1

f− (x) = f (x) − f (−x) .

2

Now f+ and f− are even and odd functions and f = f+ + f− . Thus any function in F can be

split into two components f+ and f− , such that f+ ∈ F+ and f− ∈ F− . To show that the sum

1425

T

is direct, suppose f is an element in F+ F− . Then we have that f (x) = −f (−x) = −f (x),

so f (x) = 0 for all x, i.e., f is the zero vector in F . We have established equation 364.12.1.

Version: 2 Owner: mathcam Author(s): matte

365.13 even/odd function

Definition.

Let f be a function from R to R. If f (x) = f (−x) for all x ∈ R, then f is an even function.

Similarly, if f (x) = −f (−x) for all x ∈ R, then f is an odd function.

Example.

1. The trigonometric functions sin and cos are odd and even, respectively.

properties.

**1. The vector space of real functions can be written as the direct sum of even and odd
**

functions. (See this page.)

2. Let f : R → R be a differentiable function.

**(a) If f is an even function, then the derivative f 0 is an odd function.
**

(b) If f is an odd function, then the derivative f 0 is an even function.

(proof)

**3. Let f : R → R be a smooth function. Then there exists smooth functions g, h : R → R
**

such that

f (x) = g(x2 ) + xh(x2 )

for all x ∈ R. Thus, if f is even, we have f (x) = g(x2 ), and if f is odd, we have

f (x) = xh(x2 ) ([4], Exercise 1.2)

REFERENCES

1. L. Hörmander, The Analysis of Linear Partial Differential Operators I, (Distribution

theory and Fourier Analysis), 2nd ed, Springer-Verlag, 1990.

Version: 4 Owner: mathcam Author(s): matte

1426

365.14 example of chain rule

**Suppose we wanted to differentiate
**

p

h(x) = sin(x).

Here, h(x) is given by the composition

h(x) = f (g(x)),

where √

f (x) = x and g(x) = sin(x).

Then chain rule says that

h0 (x) = f 0 (g(x))g 0(x).

Since

1

f 0 (x) = √ , and g 0 (x) = cos(x),

2 x

we have by chain rule

0 1 cos x

h (x) = √ cos x = √

2 sin x 2 sin x

**Using the Leibniz formalism, the above calculation would have the following appearance.
**

First we describe the functional relation as

p

z = sin(x).

**Next, we introduce an auxiliary variable y, and write
**

√

z= y, y = sin(x).

**We then have
**

dz 1 dy

= √ , = cos(x),

dy 2 y dx

and hence the chain rule gives

dz 1

= √ cos(x)

dx 2 y

1 cos(x)

= p

2 sin(x)

Version: 1 Owner: rmilson Author(s): rmilson

1427

365.15 example of increasing/decreasing/monotone func-

tion

**The function f (x) = ex is strictly increasing and hence strictly monotone. Similarly g(x) =
**

e−x is strictly decreasing

p and√hence strictlypmonotone. Consider the function h : [1, 10] 7→

√

[1, 5] where h(x) = x − 4 x − 1 + 3 + x − 6 x − 1 + 8. It is not strictly monotone

since it is constant on an interval, however it is decreasing and hence monotone.

Version: 1 Owner: Johan Author(s): Johan

365.16 extended mean-value theorem

**Let f : [a, b] → R and g : [a, b] → R be continuous on [a, b] and differentiable on (a, b). Then
**

there exists some number ξ ∈ (a, b) satisfying:

(f (b) − f (a))g 0 (ξ) = (g(b) − g(a))f 0(ξ).

If g is linear this becomes the usual mean-value theorem.

Version: 6 Owner: mathwizard Author(s): mathwizard

365.17 increasing/decreasing/monotone function

Definition Let A a subset of R, and let f be a function from f : A → R. Then

1. f is increasing, if x ≤ y implies that f (x) ≤ f (y) (for all x and y in A).

2. f is strictly increasing, if x < y implies that f (x) < f (y).

3. f is decreasing, if x ≥ y implies that f (x) ≥ f (y).

4. f is strictly decreasing, if x > y implies that f (x) > f (y).

5. f is monotone, if f is either increasing or decreasing.

6. f is strictly monotone, if f is either strictly increasing or strictly decreasing.

**Theorem Let X be a bounded of unbounded open S interval of R. In other words, let X be
**

an interval of the form X = (a, b), where a, b ∈ R {−∞, ∞}. Futher, let f : X → R be a

monotone function.

1428

1. The set of points where f is discontinuous is at most countable [1, 1].

Lebesgue f is differentiable almost everywhere ([1], pp. 514).

REFERENCES

1. C.D. Aliprantis, O. Burkinshaw, Principles of Real Analysis, 2nd ed., Academic Press,

1990.

2. W. Rudin, Principles of Mathematical Analysis, McGraw-Hill Inc., 1976.

3. F. Jones, Lebesgue Integration on Euclidean Spaces, Jones and Barlett Publishers, 1993.

Version: 3 Owner: matte Author(s): matte

365.18 intermediate value theorem

**Let f be a continuous function on the interval [a, b]. Let x1 and x2 be points with a ≤ x1 <
**

x2 ≤ b such that f (x1 ) 6= f (x2 ). Then for each value y between f (x1 ) and (x2 ), there is a

c ∈ (x1 , x2 ) such that f (c) = y.

Bolzano’s theorem is a special case of this one.

Version: 2 Owner: drini Author(s): drini

365.19 limit

**Let f : X \ {a} −→ Y be a function between two metric spaces X and Y , defined everywhere
**

except at some a ∈ X. For L ∈ Y , we say the limit of f (x) as x approaches a is equal to L,

or

lim f (x) = L

x→a

**if, for every real number ε > 0, there exists a real number δ > 0 such that, whenever x ∈ X
**

with 0 < dX (x, a) < δ, then dY (f (x), L) < ε.

**The formal definition of limit as given above has a well–deserved reputation for being no-
**

toriously hard for inexperienced students to master. There is no easy fix for this problem,

since the concept of a limit is inherently difficult to state precisely (and indeed wasn’t even

accomplished historically until the 1800’s by Cauchy, well after the invention of calculus in

the 1600’s by Newton and Leibniz). However, there are number of related definitions, which,

taken together, may shed some light on the nature of the concept.

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• The notion of a limit can be generalized to mappings between arbitrary topological spaces.

In this context we say that limx→a f (x) = L if and only if, for every neighborhood V

of L (in Y ), there is a deleted neighborhood U of a (in X) which is mapped into V by

f.

• Let an , n ∈ N be a sequence of elements in a metric space X. We say that L ∈ X is

the limit of the sequence, if for every ε > 0 there exists a natural number N such that

d(an , L) < ε for all natural numbers n > N.

• The definition of the limit of a mapping can be based on the limit of a sequence. To

wit, limx→a f (x) = L if and only if, for every sequence of points xn in X converging to

a (that is, xn → a, xn 6= a), the sequence of points f (xn ) in Y converges to L.

**In calculus, X and Y are frequently taken to be Euclidean spaces Rn and Rm , in which case
**

the distance functions dX and dY cited above are just Euclidean distance.

Version: 5 Owner: djao Author(s): rmilson, djao

365.20 mean value theorem

Mean value theorem Let f : [a, b] → R be a continuous function differentiable on (a, b).

**Then there is some real number x0 ∈ (a, b) such that
**

f (b) − f (a)

f 0 (x0 ) = .

b−a

Version: 3 Owner: drini Author(s): drini, apmxi

365.21 mean-value theorem

**Let f : R → R be a function which is continuous on the interval [a, b] and differentiable on
**

(a, b). Then there exists a number c : a < c < b such that

f (b) − f (a)

f 0 (c) = . (365.21.1)

b−a

The geometrical meaning of this theorem is illustrated in the picture:

1430

This is often used in the integral context: ∃c ∈ [a, b] such that

(b − a)f (c) = intba f (x)dx. (365.21.2)

Version: 4 Owner: mathwizard Author(s): mathwizard, drummond

365.22 monotonicity criterion

**Suppose that f : [a, b] → R is a function which is continuous on [a, b] and differentiable on
**

(a, b).

Then the following relations hold.

1. f 0 (x) ≥ 0 for all x ∈ (a, b) ⇔ f is an increasing function on [a, b];

2. f 0 (x) ≤ 0 for all x ∈ (a, b) ⇔ f is a decreasing function on [a, b];

3. f 0 (x) > 0 for all x ∈ (a, b) ⇒ f is a strictly increasing function on [a, b];

4. f 0 (x) < 0 for all x ∈ (a, b) ⇒ f is a strictly decreasing function on [a, b].

**Notice that the third and fourth statement cannot be inverted. As an example consider the
**

function f : [−1, 1] → R, f (x) = x3 . This is a strictly increasing function, but f 0 (0) = 0.

Version: 4 Owner: paolini Author(s): paolini

365.23 nabla

**Let f : Rn → R a C 1 (Rn ) function. That is, a partially differentiable function on all its
**

coordinates. The symbol ∇, named nabla represents the gradient operator whose action on

f (x1 , x2 , . . . , xn ) is given by

**∇f = (fx1 , fx2 , . . . , fxn )
**

∂f ∂f ∂f

= , ,...,

∂x1 ∂x2 ∂xn

Version: 2 Owner: drini Author(s): drini, apmxi

1431

365.24 one-sided limit

**Let f be a real-valued function defined on S ⊆ R. The left-hand one-sided limit at a is
**

defined to be the real number L− such that for every > 0 there exists a δ > 0 such that

|f (x) − L− | < whenever 0 < a − x < δ.

**Analogously, the right-hand one-sided limit at a is the real number L+ such that for
**

every > 0 there exists a δ > 0 such that |f (x) − L+ | < whenever 0 < x − a < δ.

Common notations for the one-sided limits are

L+ = f (x+) = lim+ f (x) = lim f (x),

x→a x&a

−

L = f (x−) = lim− f (x) = lim f (x).

x→a x%a

**Sometimes, left-handed limits are referred to as limits from below while right-handed limits
**

are from above.

**Theorem The ordinary limit of a function exists at a point if and only if both one-sided
**

limits exist at this point and are equal (to the ordinary limit).

**e.g., The Heaviside unit step function, sometimes colloquially referred to as the diving board
**

function, defined by (

0 if x < 0

H(x) =

1 if x > 0

has the simplest kind of discontinuity at x = 0, a jump discontinuity. Its ordinary limit does

not exist at this point, but the one-sided limits do exist, and are

**lim H(x) = 0 and lim+ H(x) = 1.
**

x→0− x→0

Version: 5 Owner: matte Author(s): matte, NeuRet

365.25 product rule

**The product rule states that if f : R → R and g : R → R are functions in one variable
**

both differentiable at a point x0 , then the derivative of the product of the two fucntions,

denoted f · g, at x0 is given by

D

(f · g) (x0 ) = f (x0 )g 0 (x0 ) + f 0 (x0 )g(x0 ).

Dx

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Proof

See the proof of the product rule.

365.25.1 Generalized Product Rule

**More generally, for differentiable functions f1 , f2 , . . . , fn in one variable, all differentiable at
**

x0 , we have

n

X

D(f1 · · · fn )(x0 ) = (fi (x0 ) · · · fi−1 (x0 ) · Dfi (x0 ) · fi+1 (x0 ) · · · fn (x0 )) .

i=1

Also see Leibniz’ rule.

Example

**The derivative of x ln |x| can be found by application of this rule. Let f (x) = x, g(x) = ln |x|,
**

so that f (x)g(x) = x ln |x|. Then f 0 (x) = 1 and g 0 (x) = x1 . Therefore, by the product rule,

D

(x ln |x|) = f (x)g 0 (x) + f 0 (x)g(x)

Dx

x

= + 1 · ln |x|

x

= ln |x| + 1

Version: 8 Owner: mathcam Author(s): mathcam, Logan

365.26 proof of Darboux’s theorem

**WLOG, assume f+0 (a) > t > f−0 (b). Let g(x) = f (x) − tx. Then g 0 (x) = f 0 (x) − t,
**

0 0

g+ (a) > 0 > g− (b), and we wish to find a zero of g 0 .

**g is a continuous function on [a, b], so it attains a maximum on [a, b]. This maximum cannot
**

0 0

be at a, since g+ (a) > 0 so g is locally increasing at a. Similarly, g− (b) < 0, so g is locally

decreasing at b and cannot have a maximum at b. So the maximum is attained at some

c ∈ (a, b). But then g 0 (c) = 0 by Fermat’s theorem.

Version: 2 Owner: paolini Author(s): paolini, ariels

1433

365.27 proof of Fermat’s Theorem (stationary points)

**Suppose that x0 is a local maximum (a similar proof applies if x0 is a local minimum). Then
**

there exists δ > 0 such that (x0 − δ, x0 + δ) ⊂ (a, b) and such that we have f (x0 ) ≥ f (x) for

all x with |x − x0 | < δ. Hence for h ∈ (0, δ) we notice that it holds

f (x0 + h) − f (x0 )

≤ 0.

h

Since the limit of this ratio as h → 0+ exists and is equal to f 0 (x0 ) we conclude that

f 0 (x0 ) ≤ 0. On the other hand for h ∈ (−δ, 0) we notice that

f (x0 + h) − f (x0 )

≥0

h

but again the limit as h → 0+ exists and is equal to f 0 (x0 ) so we also have f 0 (x0 ) ≥ 0.

Hence we conclude that f 0 (x0 ) = 0.

Version: 1 Owner: paolini Author(s): paolini

365.28 proof of Rolle’s theorem

**Because f is continuous on a compact (closed and bounded) interval I = [a, b], it attains its
**

maximum and minimum values. In case f (a) = f (b) is both the maximum and the minimum,

then there is nothing more to say, for then f is a constant function and f 0 ⇔ 0 on the whole

interval I. So suppose otherwise, and f attains an extremum in the open interval (a, b), and

without loss of generality, let this extremum be a maximum, considering −f in lieu of f as

necessary. We claim that at this extremum f (c) we have f 0 (c) = 0, with a < c < b.

**To show this, note that f (x) − f (c) 6 0 for all x ∈ I, because f (c) is the maximum. By
**

definition of the derivative, we have that

f (x) − f (c)

f 0 (c) = lim .

x→c x−c

Looking at the one-sided limits, we note that

f (x) − f (c)

R = lim+ 60

x→c x−c

because the numerator in the limit is nonpositive in the interval I, yet x − c > 0, as x

approaches c from the right. Similarly,

f (x) − f (c)

L = lim− > 0.

x→c x−c

Since f is differentiable at c, the left and right limits must coincide, so 0 6 L = R 6 0, that

is to say, f 0 (c) = 0.

Version: 1 Owner: rmilson Author(s): NeuRet

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365.29 proof of Taylor’s Theorem

**Let n be a natural number and I be the closed interval [a, b]. We have that f : I → R has
**

n continuous derivatives and its (n + 1)-st derivative exists. Suppose that c ∈ I, and x ∈ I

is arbitrary. Let J be the closed interval with endpoints c and x.

Define F : J → R by

n

X (x − t)k

F (t) := f (x) − f (k) (t) (365.29.1)

k=0

k!

so that

n

X

0 0 (x − t)k (k+1) (x − t)k−1 (k)

F (t) = f (t) − f (t) − f (t)

k=1

k! (k − 1)!

n

(x − t) (n+1)

= − f (t)

n!

since the sum telescopes. Now, define G on J by

n+1

x−t

G(t) := F (t) − F (c)

x−c

**and notice that G(c) = G(x) = 0. Hence, Rolle’s theorem gives us a ζ strictly between x
**

and c such that

(x − ζ)n

0 = G0 (ζ) = F 0 (ζ) − (n + 1) F (c)

(x − c)n+1

that yields

1 (x − c)n+1 0

F (c) = − F (ζ)

n + 1 (x − c)n

1 (x − c)n+1 (x − ζ)n (n+1)

= f (ζ)

n + 1 (x − c)n n!

f (n+1) (ζ)

= (x − c)n+1

(n + 1)!

**from which we conclude, recalling (364.29.1),
**

n

X f (k) (c) f (n+1) (ζ)

f (x) = (x − c)k + (x − c)n+1

k=0

k! (n + 1)!

Version: 3 Owner: rmilson Author(s): NeuRet

1435

365.30 proof of binomial formula

**Let p ∈ R and x ∈ R, |x| < 1 be given. We wish to show that
**

X∞

p xn

(1 + x) = pn ,

n=0

n!

where pn denotes the nth falling factorial of p.

**The convergence of the series in the right-hand side of the above equation is a straight-
**

forward consequence of the ratio test. Set

f (x) = (1 + x)p .

and note that

f (n) (x) = pn (1 + x)p−n .

The desired equality now follows from Taylor’s Theorem. Q.E.D.

Version: 2 Owner: rmilson Author(s): rmilson

365.31 proof of chain rule

**Let’s say that g is differentiable in x0 and f is differentiable in y0 = g(x0 ). We define:
**

f (y)−f (y0 )

y−y0

if y 6= y0

ϕ(y) = 0

f (y0 ) if y = y0

**Since f is differentiable in y0 , ϕ is continuous. We observe that, for x 6= x0 ,
**

f (g(x)) − f (g(x0 )) g(x) − g(x0 )

= ϕ(g(x)) ,

x − x0 x − x0

in fact, if g(x) 6= g(x0 ), it follows at once from the definition of ϕ, while if g(x) = g(x0 ),

both members of the equation are 0.

**Since g is continuous in x0 , and ϕ is continuous in y0 ,
**

lim ϕ(g(x)) = ϕ(g(x0 )) = f 0 (g(x0 )),

x→x0

hence

f (g(x)) − f (g(x0 ))

(f ◦ g)0(x0 ) = lim

x→x0 x − x0

g(x) − g(x0 )

= lim ϕ(g(x))

x→x0 x − x0

0 0

= f (g(x0 ))g (x0 ).

Version: 3 Owner: n3o Author(s): n3o

1436

365.32 proof of extended mean-value theorem

**Let f : [a, b] → R and g : [a, b] → R be continuous on [a, b] and differentiable on (a, b).
**

Define the function

h(x) = f (x) (g(b) − g(a)) − g(x) (f (b) − f (a)) − f (a)g(b) + f (b)g(a).

**Because f and g are continuous on [a, b] and differentiable on (a, b), so is h. Furthermore,
**

h(a) = h(b) = 0, so by Rolle’s theorem there exists a ξ ∈ (a, b) such that h0 (ξ) = 0. This

implies that

f 0 (ξ) (g(b) − g(a)) − g 0(ξ) (f (b) − f (a)) = 0

and, if g(b) 6= g(a),

f 0 (ξ) f (b) − f (a)

0

= .

g (ξ) g(b) − g(a)

Version: 3 Owner: pbruin Author(s): pbruin

365.33 proof of intermediate value theorem

We first prove the following lemma.

**If f : [a, b] → R is a continuous function with f (a) ≤ 0 ≤ f (b) then ∃c ∈ [a, b] such that
**

f (c) = 0.

Define the sequences (an ) and (bn ) inductively, as follows.

a0 = a b0 = b

an + bn

cn =

2

(

(an−1 , cn−1 ) f (cn−1 ) ≥ 0

(an , bn ) =

(cn−1 , bn−1 ) f (cn−1 ) < 0

We note that

a0 ≤ a1 . . . ≤ an ≤ bn ≤ . . . b1 ≤ b0

(bn − an ) = 2−n (b0 − a0 ) (365.33.1)

f (an ) ≤ 0 ≤ f (bn ) (365.33.2)

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By the fundamental axiom of analysis (an ) → α and (bn ) → β. But (bn − an ) → 0 so α = β.

By continuity of f

(f (an )) → f (α) (f (bn )) → f (α)

But we have f (α) ≤ 0 and f (α) ≥ 0 so that f (α) = 0. Furthermore we have a ≤ α ≤ b,

proving the assertion.

**Set g(x) = f (x) − k where f (a) ≤ k ≤ f (b). g satisfies the same conditions as before, so ∃c
**

such that f (c) = k. Thus proving the more general result.

Version: 2 Owner: vitriol Author(s): vitriol

365.34 proof of mean value theorem

Define h(x) on [a, b] by

f (b) − f (a)

h(x) = f (x) − f (a) − (x − a)

b−a

clearly, h is continuous on [a, b], differentiable on (a, b), and

h(a) = f (a) − f (a) = 0

f (b)−f (a)

h(b) = f (b) − f (a) − b−a

(b − a) = 0

**Notice that h satisfies the conditions of Rolle’s theorem. Therefore, by Rolle’s Theorem
**

there exists c ∈ (a, b) such that h0 (c) = 0.

**However, from the definition of h we obtain by differentiation that
**

f (b) − f (a)

h0 (x) = f 0 (x) −

b−a

Since h0 (c) = 0, we therefore have

f (b) − f (a)

f 0 (c) =

b−a

as required.

REFERENCES

1. Michael Spivak, Calculus, 3rd ed., Publish or Perish Inc., 1994.

Version: 2 Owner: saforres Author(s): saforres

1438

365.35 proof of monotonicity criterion

Let us start from the implications “⇒”.

**Suppose that f 0 (x) ≥ 0 for all x ∈ (a, b). We want to prove that therefore f is increasing. So
**

take x1 , x2 ∈ [a, b] with x1 < x2 . Applying the mean-value theorem on the interval [x1 , x2 ]

we know that there exists a point x ∈ (x1 , x2 ) such that

f (x2 ) − f (x1 ) = f 0 (x)(x2 − x1 )

and being f 0 (x) ≥ 0 we conclude that f (x2 ) ≥ f (x1 ).

This proves the first claim. The other three cases can be achieved with minor modifications:

replace all “≥” respectively with ≤, > and <.

Let us now prove the implication “⇐” for the first and second statement.

**Given x ∈ (a, b) consider the ratio
**

f (x + h) − f (x)

.

h

If f is increasing the numerator of this ratio is ≥ 0 when h > 0 and is ≤ 0 when h < 0.

Anyway the ratio is ≥ 0 since the denominator has the same sign of the numerator. Since

we know by hypothesys that the function f is differentiable in x we can pass to the limit to

conclude that

f (x + h) − f (x)

f 0 (x) = lim ≥ 0.

h→0 h

If f is decreasing the ratio considered turns out to be ≤ 0 hence the conclusion f 0 (x) ≤ 0.

**Notice that if we suppose that f is strictly increasing we obtain the this ratio is > 0, but
**

passing to the limit as h → 0 we cannot conclude that f 0 (x) > 0 but only (again) f 0 (x) ≥ 0.

Version: 2 Owner: paolini Author(s): paolini

365.36 proof of quotient rule

Let F (x) = f (x)/g(x). Then

f (x+h) f (x)

0 F (x + h) − F (x) g(x+h)

− g(x)

F (x) = lim = lim h

h→0 h h→0 h

f (x + h)g(x) − f (x)g(x + h)

= lim

h→0 hg(x + h)g(x)

1439

Like the product rule, the key to this proof is subtracting and adding the same quantity. We

separate f and g in the above expression by subtracting and adding the term f (x)g(x) in

the numerator.

f (x + h)g(x) − f (x)g(x) + f (x)g(x) − f (x)g(x + h)

F 0 (x) = lim

h→0 hg(x + h)g(x)

g(x) f (x+h)−f

h

(x)

− f (x) g(x+h)−g(x)

h

= lim

h→0 g(x + h)g(x)

limh→0 g(x) · limh→0 f (x+h)−f

h

(x)

− limh→0 f (x) · limh→0 g(x+h)−g(x)

h

=

limh→0 g(x + h) · limh→0 g(x)

g(x)f 0(x) − f (x)g 0 (x)

=

[g(x)]2

Version: 1 Owner: Luci Author(s): Luci

365.37 quotient rule

**The quotient rule says that the derivative of the quotient f /g of two differentiable functions
**

f and g exists at all values of x as long as g(x) 6= 0 and is given by the formula

d f (x) g(x)f 0 (x) − f (x)g 0(x)

=

dx g(x) [g(x)]2

**The Quotient Rule and the other differentiation formulas allow us to compute the derivative
**

of any rational function.

Version: 10 Owner: Luci Author(s): Luci

365.38 signum function

**The signum function is the function sign : R → R
**

−1 when x < 0,

sign(x) = 0 when x = 0,

1 when x > 0.

The following properties hold:

1440

1. For all x ∈ R, sign(−x) = − sign(x).

**2. For all x ∈ R, |x| = sign(x)x.
**

d

3. For all x 6= 0, dx

|x| = sign(x).

**Here, we should point out that the signum function is often defined simply as 1 for x > 0 and
**

−1 for x < 0. Thus, at x = 0, it is left undefined. See e.g. [2]. In applications, such as the

Laplace transform, this definition is adequate since the value of a function at a single point

does not change the analysis. One could then, in fact, set sign(0) to any value. However,

setting sign(0) = 0 is motivated by the above relations.

**A related function is the Heaviside step function defined as
**

0 when x < 0,

H(x) = 1/2 when x = 0,

1 when x > 0.

**Again, this function is sometimes left undefined at x = 0. The motivation for setting
**

H(0) = 1/2 is that for all x ∈ R, we then have the relations

1

H(x) = (sign(x) + 1),

2

H(−x) = 1 − H(x).

**This first relation is clear. For the second, we have
**

1

1 − H(x) = 1 − (sign(x) + 1)

2

1

= (1 − sign(x))

2

1

= (1 + sign(−x))

2

= H(−x).

**Example Let a < b be real numbers, and let f : R → R be the piecewise defined function
**

4 when x ∈ (a, b),

f (x) =

0 otherwise.

**Using the Heaviside step function, we can write
**

f (x) = 4 H(x − a) − H(x − b) (365.38.1)

**almost everywhere. Indeed, if we calculate f using equation 364.38.1 we obtain f (x) = 4 for
**

x ∈ (a, b), f (x) = 0 for x ∈

/ [a, b], and f (a) = f (b) = 2. Therefore, equation 364.38.1 holds

at all points except a and b. 2

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365.38.1 Signum function for complex arguments

**For a complex number z, the signum function is defined as [1]
**

0 when z = 0,

sign(z) =

z/|z| when z = 6 0.

**In other words, if z is non-zero, then sign z is the projection of z onto the unit circle {z ∈
**

C | |z| = 1}. clearly, the complex signum function reduces to the real signum function for

real arguments. For all z ∈ C, we have

z sign z = |z|,

where z is the complex conjugate of z.

REFERENCES

1. E. Kreyszig, Advanced Engineering Mathematics, John Wiley & Sons, 1993, 7th ed.

2. G. Bachman, L. Narici, Functional analysis, Academic Press, 1966.

Version: 4 Owner: mathcam Author(s): matte

1442

Chapter 366

26A09 – Elementary functions

366.1 definitions in trigonometry

Informal definitions

**Given a triangle ABC with a signed angle x at A and a right angle at B, the ratios
**

BC AB BC

AC AC AB

are dependant only on the angle x, and therefore define functions, denoted by

sin x cos x tan x

respectively, where the names are short for sine, cosine and tangent. Their inverses are

rather less important, but also have names:

1

cot x = AB/BC = (cotangent)

tan x

1

csc x = AC/BC = (cosecant)

sin x

1

sec x = AC/AB = (secant)

cos x

From Pythagoras’s theorem we have cos2 x + sin2 x = 1 for all (real) x. Also it is “clear”

from the diagram at left that functions cos and sin are periodic with period 2π. However:

Formal definitions

The above definitions are not fully rigorous, because we have not defined the word angle.

We will sketch a more rigorous approach.

**The power series
**

∞

X xn

n=0

n!

1443

converges uniformly on compact subsets of C and its sum, denoted by exp(x) or by ex , is

therefore an entire function of x, called the exponential function. f (x) = exp(x) is the

unique solution of the boundary value problem

f (0) = 1 f 0 (x) = f (x)

on R. The sine and cosine functions, for real arguments, are defined in terms of exp, simply

by

exp(ix) = cos x + i(sin x) .

Thus

x2 x4 x6

cos x = 1 − + − + ...

2! 4! 6!

x x3 x5

sin x = − + − ...

1! 3! 5!

Although it is not self-evident, cos and sin are periodic functions on the real line, and have

the same period. That period is denoted by 2π.

Version: 3 Owner: Daume Author(s): Larry Hammick

366.2 hyperbolic functions

The hyperbolic functions sinh x and cosh x ared defined as follows:

ex − e−x

sinh x :=

2

e + e−x

x

cosh x := .

2

One can then also define the functions tanh x and coth x in analogy to the definitions of

tan x and cot x:

sinh x ex − e−x

tanh x := = x

cosh x e + e−x

coth x ex + e−x

coth x := = x .

cosh x e − e−x

The hyperbolic functions are named in that way because the hyperbola

x2 y 2

− 2 =1

a2 b

can be written in parametrical form with the equations:

x = a cosh t, y = b sinh t.

1444

This is because of the equation

cosh2 x − sinh2 x = 1.

There are also addition formulas which are like the ones for trigonometric functions:

**sinh(x ± y) = sinh x cosh y ± cosh x sinh y
**

cosh(x ± y) = cosh x cosh y ± sinh x sinh y.

**The Taylor series for the hyperbolic functions are:
**

∞

X x2n+1

sinh x =

n=0

(2n + 1)!

X∞

x2n

cosh x = .

n=0

(2n)!

**Using complex numbers we can use the hyperbolic functions to express the trigonometric
**

functions:

sinh(ix)

sin x =

i

cos x = cosh(ix).

Version: 2 Owner: mathwizard Author(s): mathwizard

1445

Chapter 367

**26A12 – Rate of growth of functions,
**

orders of infinity, slowly varying

functions

367.1 Landau notation

Given two functions f and g from R+ to R+ , the notation

f = O(g)

f (x)

means that the ratio g(x)

stays bounded as x → ∞. If moreover that ratio approaches zero,

we write

f = o(g).

**It is legitimate to write, say, 2x = O(x) = O(x2 ), with the understanding that we are using
**

the equality sign in an unsymmetric (and informal) way, in that we do not have, for example,

O(x2 ) = O(x).

The notation

f = Ω(g)

f (x)

means that the ratio g(x)

is bounded away from zero as x → ∞, or equivalently g = O(f ).

If both f = O(g) and f = Ω(g), we write f = Θ(g).

One more notational convention in this group is

f (x) ∼ g(x),

meaning limx→∞ fg(x)

(x)

= 1.

1446

In analysis, such notation is useful in describing error estimates. For example, the Riemann hypothesis

is equivalent to the conjecture

x √

π(x) = + O( x log x)

log x

**Landau notation is also handy in applied mathematics, e.g. in describing the efficiency of an
**

algorithm. It is common to say that an algorithm requires O(x3 ) steps, for example, without

needing to specify exactly what is a step; for if f = O(x3 ), then f = O(Ax3 ) for any positive

constant A.

Version: 8 Owner: mathcam Author(s): Larry Hammick, Logan

1447

Chapter 368

**26A15 – Continuity and related
**

questions (modulus of continuity,

semicontinuity, discontinuities, etc.)

368.1 Dirichlet’s function

**Dirichlet’s function f : R → R is defined as
**

1

q

if x = pq is a rational number in lowest terms,

f (x) =

0 if x is an irrational number.

**This function has the property that it is continuous at every irrational number and discontinuous
**

at every rational one.

Version: 3 Owner: urz Author(s): urz

368.2 semi-continuous

A real function f : A → R, where A ⊆ R is said to be lower semi-continuous in x0 if

∀ε > 0 ∃δ > 0 ∀x ∈ A |x − x0 | < δ ⇒ f (x) > f (x0 ) − ε,

and f is said to be upper semi-continuous if

∀ε > 0 ∃δ > 0 ∀x ∈ A |x − x0 | < δ ⇒ f (x) < f (x0 ) + ε.

1448

Remark A real function is continuous in x0 if and only if it is both upper and lower semicontinuous

in x0 .

We can generalize the definition to arbitrary topological spaces as follows.

**Let A be a topological space. f : A → R is lower semicontinuous at x0 if, for each ε > 0
**

there is a neighborhood U of x0 such that x ∈ U implies f (x) > f (x0 ) − ε.

**Theorem Let f : [a, b] → R be a lower (upper) semi-continuous function. Then f has a
**

minimum (maximum) in [a, b].

Version: 3 Owner: drini Author(s): drini, n3o

368.3 semicontinuous

**Defintion [1] Suppose X is a topological space, and f is a function from X into the
**

extended real numbers R; f : X → R. Then:

**1. If {x ∈ X | f (x) > α} is an open set in X for all α ∈ R, then f is said to be lower
**

semicontinuous.

**2. If {x ∈ X | f (x) < α} is an open set in X for all α ∈ R, then f is said to be upper
**

semicontinuous.

Properties

**1. If X is a topological space and f is a function f : X → R, then f is continuous if and
**

only if f is upper and lower semicontinuous [1, 3].

2. The characteristic function of an open set is lower semicontinuous [1, 3].

3. The characteristic function of a closed set is upper semicontinuous [1, 3].

4. If f and g are lower semicontinuous, then f + g is also lower semicontinuous [3].

REFERENCES

1. W. Rudin, Real and complex analysis, 3rd ed., McGraw-Hill Inc., 1987.

2. D.L. Cohn, Measure Theory, Birkhäuser, 1980.

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1449

368.4 uniformly continuous

**Let f : A → R be a real function defined on a subset A of the real line. We say that f is
**

uniformly continuous if, given an arbitrary small positive ε, there exists a positive δ such

that whenever two points in A differ by less than δ, they are mapped by f into points which

differ by less than ε. In symbols:

∀ε > 0 ∃δ > 0 ∀x, y ∈ A |x − y| < δ ⇒ |f (x) − f (y)| < ε.

**Every uniformly continuous function is also continuous, while the converse does not always
**

hold. For instance, the function f :]0, +∞[→ R defined by f (x) = 1/x is continuous in its

domain, but not uniformly.

**A more general definition of uniform continuity applies to functions between metric spaces
**

(there are even more general environments for uniformly continuous functions, i.e. Uniform spaces).

Given a function f : X → Y , where X and Y are metric spaces with distances dX and dY ,

we say that f is uniformly continuous if

∀ε > 0 ∃δ > 0 ∀x, y ∈ X dX (x, y) < δ ⇒ dY (f (x), f (y)) < ε.

**Uniformly continuous functions have the property that they map Cauchy sequences to Cauchy
**

sequences and that they preserve uniform convergence of sequences of functions.

Any continuous function defined on a compact space is uniformly continuous (see Heine-Cantor theorem).

Version: 10 Owner: n3o Author(s): n3o

1450

Chapter 369

26A16 – Lipschitz (Hölder) classes

369.1 Lipschitz condition

**A mapping f : X → Y between metric spaces is said to satisfy the Lipschitz condition if
**

there exists a real constant α > 0 such that

dY (f (p), f (q)) 6 αdX (p, q), for all p, q ∈ X.

Proposition 17. A Lipschitz mapping f : X → Y is uniformly continuous.

**Proof. Let f be a Lipschitz mapping and α > 0 a corresponding Lipschitz constant. For
**

every given > 0, choose δ > 0 such that

δα < .

**Let p, q ∈ X such that
**

dX (p, q) < δ

be given. By assumption,

dY (f (p), f (q)) 6 αδ < ,

as desired. QED

**Notes. More generally, one says that mapping satisfies a Lipschitz condition of order β > 0
**

if there exists a real constant α > 0 such that

dY (f (p), f (q)) 6 αdX (p, q)β , for all p, q ∈ X.

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1451

369.2 Lipschitz condition and differentiability

**If X and Y are Banach spaces, e.g. Rn , one can inquire about the relation between differ-
**

entiability and the Lipschitz condition. The latter is the weaker condition. If f is Lipschitz,

the ratio

kf (q) − f (p)k

, p, q ∈ X

kq − pk

is bounded but is not assumed to converge to a limit. Indeed, differentiability is the stronger

condition.

**Proposition 18. Let f : X → Y be a continuously differentiable mapping between Banach
**

spaces. If K ⊂ X is a compact subset, then the restriction f : K → Y satisfies the Lipschitz

condition.

**Proof. Let Lin(X, Y ) denote the Banach space of bounded linear maps from X to Y . Recall
**

that the norm kT k of a linear mapping T ∈ Lin(X, Y ) is defined by

kT uk

kT k = sup{ : u 6= 0}.

kuk

**Let Df : X → Lin(X, Y ) denote the derivative of f . By definition Df is continuous, which
**

really means that kDf k : X → R is a continuous function. Since K ⊂ X is compact, there

exists a finite upper bound B1 > 0 for kDf k restricted to U. In particular, this means that

kDf (p)uk 6 kDf (p)kkuk 6 B1 kuk,

for all p ∈ K, u ∈ X.

**Next, consider the secant mapping s : X × X → R defined by
**

kf (q) − f (p) − Df (p)(q − p)k q 6= p

s(p, q) = kq − pk

0 p=q

**This mapping is continuous, because f is assumed to be continuously differentiable. Hence,
**

there is a finite upper bound B2 > 0 for s restricted to the compact K × K. It follows that

for all p, q ∈ K we have

kf (q) − f (p)k 6 kf (q) − f (p) − Df (p)(q − p)k + kDf (p)(q − p)k

6 B2 kq − pk + B1 kq − pk

= (B1 + B2 )kq − pk

Therefore B1 , B2 is the desired Lipschitz constant. QED

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369.3 Lipschitz condition and differentiability result

About Lipschitz continuity of differentiable functions the following holds.

**Theorem 6. Let X, Y be Banach spaces and let A be a convex (see convex set), open subset
**

of X. Let f : A → Y be a function which is continuous in A and differentiable in A. Then

f is lipschitz continuous on A if and only if the derivative Df is bounded on A i.e.

sup kDf (x)k < +∞.

x∈A

S uppose that f is lipschitz continuous:

kf (x) − f (y)k ≤ Lkx − yk.

Then given any x ∈ A and any v ∈ X, for all small h ∈ R we have

f (x + hv) − f (x)

k k ≤ L.

h

Hence, passing to the limit h → 0 it must hold kDf (x)k ≤ L.

On the other hand suppose that Df is bounded on A:

kDf (x)k ≤ L, ∀x ∈ A.

Given any two points x, y ∈ A and given any α ∈ Y ∗ consider the function G : [0, 1] → R

G(t) = hα, f ((1 − t)x + ty)i.

For t ∈ (0, 1) it holds

G0 (t) = hα, Df ((1 − t)x + ty)[y − x]i

and hence

|G0 (t)| ≤ Lkαk ky − xk.

Applying Lagrange mean-value theorem to G we know that there exists ξ ∈ (0, 1) such that

|hα, f (y) − f (x)i| = |G(1) − G(0)| = |G0(ξ)| ≤ kαkLky − xk

and since this is true for all α ∈ Y ∗ we get

kf (y) − f (x)k ≤ Lky − xk

which is the desired claim.

Version: 1 Owner: paolini Author(s): paolini

1453

Chapter 370

26A18 – Iteration

370.1 iteration

**Let f : X → X be a function, X being any set. The n-th iteration of a function is the
**

function which is obtained if f is applied n times, and is denoted by f n . More formally we

define:

f 0 (x) = x

and

f n+1 (x) = f (f n (x))

for nonnegative integers n. If f is invertible, then by going backwards we can define the

iterate also for negative n.

Version: 6 Owner: mathwizard Author(s): mathwizard

370.2 periodic point

**Let f : X → X be a function and f n its n-th iteration. A point x is called a periodic point
**

of period n of f if it is a fixed point of f n . The least n for which x is a fixed point of f n is

called prime period or least period.

**If f is a function mapping R to R or C to C then a periodic point x of prime period n is
**

called hyperbolic if |(f n )0 (x)| =

6 1, attractive if |(f n )0 (x)| < 1 and repelling if |(f n )0 (x)| > 1.

Version: 11 Owner: mathwizard Author(s): mathwizard

1454

Chapter 371

**26A24 – Differentiation (functions of
**

one variable): general theory,

generalized derivatives, mean-value

theorems

371.1 Leibniz notation

**Leibniz notation centers around the concept of a differential element. The differential
**

element of x is represented by dx. You might think of dx as being an infinitesimal change

dy

in x. It is important to note that d is an operator, not a variable. So, when you see dx , you

y

can’t automatically write as a replacement x .

df (x) d

We use dx

or dx

f (x) to represent the derivative of a function f (x) with respect to x.

df (x) f (x + Dx) − f (x)

= lim

dx Dx→0 Dx

We are dividing two numbers infinitely close to 0, and arriving at a finite answer. D is

another operator that can be thought of just a change in x. When we take the limit of Dx

as Dx approaches 0, we get an infinitesimal change dx.

**Leibniz notation shows a wonderful use in the following example:
**

dy dy du dy du

= =

dx dx du du dx

The two dus can be cancelled out to arrive at the original derivative. This is the Leibniz

notation for the chain rule.

**Leibniz notation shows up in the most common way of representing an integral,
**

F (x) = intf (x)dx

1455

The dx is in fact a differential element. Let’s start with a derivative that we know (since

F (x) is an antiderivative of f (x)).

dF (x)

= f (x)

dx

dF (x) = f (x)dx

intdF (x) = intf (x)dx

F (x) = intf (x)dx

We can think of dF (x) as the differential element of area. Since dF (x) = f (x)dx, the element

of area is a rectangle, with f (x) × dx as its dimensions. Integration is the sum of all these

infinitely thin elements of area along a certain interval. The result: a finite number.

(a diagram is deserved here)

**One clear advantage of this notation is seen when finding the length s of a curve. The
**

formula is often seen as the following:

s = intds

The length is the sum of all the elements,qds, of length. If we have a function f (x), the

length element is usually written as ds = 1 + [ dfdx (x) 2

] dx. If we modify this a bit, we get

p

2 2

ds = [dx] + [df (x)] . Graphically, we could say that the length element is the hypotenuse

of a right triangle with one leg being the x element, and the other leg being the f (x) element.

(another diagram would be nice!)

**There are a few caveats, such as if you want to take the value of a derivative. Compare to
**

the prime notation.

0 df (x)

f (a) =

dx x=a

**A second derivative is represented as follows:
**

d dy d2 y

= 2

dx dx dx

3

d y

The other derivatives follow as can be expected: dx 3 , etc. You might think this is a little

**sneaky, but it is the notation. Properly using these terms can be interesting. For example,
**

2 d2 y

what is int ddxy ? We could turn it into int dx dy dy

2 dx or intd dx . Either way, we get dx .

Version: 2 Owner: xriso Author(s): xriso

371.2 derivative

**Qualitatively the derivative is a measure of the change of a function in a small region around
**

a specified point.

1456

Motivation

**The idea behind the derivative comes from the straight line. What characterizes a straight
**

line is the fact that it has constant “slope”.

Figure 371.1: The straight line y = mx + b

**In other words for a line given by the equation y = mx + b, as in Fig. 370.1, the ratio of ∆y
**

∆y

over ∆x is always constant and has the value ∆x = m.

Figure 371.2: The parabola y = x2 and its tangent at (x0 , y0 )

For other curves we cannot define a “slope”, like for the straight line, since such a quantity

would not be constant. However, for sufficiently smooth curves, each point on a curve has a

tangent line. For example consider the curve y = x2 , as in Fig. 370.2. At the point (x0 , y0 )

on the curve, we can draw a tangent of slope m given by the equation y − y0 = m(x − x0 ).

**Suppose we have a curve of the form y = f (x), and at the point (x0 , f (x0 )) we have a tangent
**

given by y − y0 = m(x − x0 ). Note that for values of x sufficiently close to x0 we can make

the approximation f (x) ≈ m(x − x0 ) + y0. So the slope m of the tangent describes how much

f (x) changes in the vicinity of x0 . It is the slope of the tangent that will be associated with

the derivative of the function f (x).

Formal definition

**More formally for any real function f : R → R, we define the derivative of f at the point x
**

as the following limit (if it exists)

f (x + h) − f (x)

f 0 (x) := lim .

h→0 h

This definition turns out to be consistent with the motivation introduced above.

The derivatives for some elementary functions are (cf. Derivative notation)

d

1. dx

c = 0, where c is constant;

d n

2. dx

x = nxn−1 ;

d

3. dx

sin x = cos x;

1457

d

4. dx

cos x = − sin x;

d x

5. dx

e = ex ;

d

6. dx

ln x = x1 .

**While derivatives of more complicated expressions can be calculated algorithmically using
**

the following rules

d

Linearity dx

(af (x) + bg(x)) = af 0 (x) + bg 0 (x);

d

Product rule dx

(f (x)g(x)) = f 0 (x)g(x) + f (x)g 0 (x);

d

Chain rule dx

g(f (x)) = g 0(f (x))f 0 (x);

d f (x) f 0 (x)g(x)−f (x)g 0 (x)

Quotient Rule dx g(x)

= g(x)2

.

**Note that the quotient rule, although given as much importance as the other rules in ele-
**

mentary calculus, can be derived by succesively applying the product rule and the chain rule

to fg(x)

(x) 1

= f (x) g(x) . Also the quotient rule does not generalize as well as the other ones.

**Since the derivative f 0 (x) of f (x) is also a function x, higher derivatives can be obtained by
**

applying the same procedure to f 0 (x) and so on.

Generalization

Banach Spaces

**Unfortunately the notion of the “slope of the tangent” does not directly generalize to more
**

abstract situations. What we can do is keep in mind the facts that the tangent is a linear

function and that it approximates the function near the point of tangency, as well as the

formal definition above.

**Very general conditions under which we can define a derivative in a manner much similar to
**

the above areas follows. Let f : V → W, where V and W are Banach spaces. Suppose that

h ∈ V and h 6= 0, the we define the directional derivative (Dh f )(x) at x as the following

limit

f (x + h) − f (x)

(Dh f )(x) := lim ,

→0

where is a scalar. Note that f (x + h) ≈ f (x) + (Dh f )(x), which is consistent with our

original motivation. This directional derivative is also called the Gâteaux derivative.

1458

Finally we define the derivative at x as the bounded linear map (Df )(x) : V → W such that

for any non-zero h ∈ V

(f (x + h) − f (x)) − (Df )(x) · h

lim = 0.

khk→0 khk

**Once again we have f (x + h) ≈ f (x) + (Df )(x) · h. In fact, if the derivative (Df )(x)
**

exists, the directional derivatives can be obtained as (Dh f )(x) = (Df )(x) · h.1 each non-

zero h ∈ V does not guarantee the existence of (Df )(x). This derivative is also called the

Fréchet derivative. In the more familiar case f : Rn → Rm , the derivative Df is simply

the Jacobian of f .

Under these general conditions the following properties of the derivative remain

1. Dh = 0, where h is a constant;

2. D(A · x) = A, where A is linear.

Linearity D(af (x) + bg(x)) · h = a(Df )(x) · h + b(Dg)(x) · h;

**“Product” rule D(B(f (x), g(x)))·h = B((Df )(x)·h, g(x))+B(f (x), (Dg)(x)·h), where
**

B is bilinear;

Chain rule D(g(f (x)) · h = (Dg)(f (x)) · ((Df )(x) · h).

**Note that the derivative of f can be seen as a function Df : V → L(V, W) given by Df : x 7→
**

(Df )(x), where L(V, W) is the space of bounded linear maps from V to W. Since L(V, W)

can be considered a Banach space itself with the norm taken as the operator norm, higher

derivatives can be obtained by applying the same procedure to Df and so on.

Manifolds

**A manifold is a topological space that is locally homeomorphic to a Banach space V (for
**

finite dimensional manifolds V = Rn ) and is endowed with enough structure to define deriva-

tives. Since the notion of a manifold was constructed specifically to generalize the notion of

a derivative, this seems like the end of the road for this entry. The following discussion is

rather technical, a more intuitive explanation of the same concept can be found in the entry

on related rates.

**Consider manifolds V and W modeled on Banach spaces V and W, respectively. Say we
**

have y = f (x) for some x ∈ V and y ∈ W , then, by definition of a manifold, we can find

1

The notation A · h is used when h is a vector and A a linear operator. This notation can be considered

advantageous to the usual notation A(h), since the latter is rather bulky and the former incorporates the

intuitive distributive properties of linear operators also associated with usual multiplication.

1459

charts (X, x) and (Y, y), where X and Y are neighborhoods of x and y, respectively. These

charts provide us with canonical isomorphisms between the Banach spaces V and W, and

the respective tangent spaces Tx V and Ty W :

dxx : Tx V → V, dyy : Ty W → W.

**Now consider a map f : V → W between the manifolds. By composing it with the chart
**

maps we construct the map

(Y,y)

g(X,x) = y ◦ f ◦ x−1 : V → W,

**defined on an appropriately restricted domain. Since we now have a map between Banach
**

(Y,y)

spaces, we can define its derivative at x(x) in the sense defined above, namely Dg(X,x) (x(x)).

If this derivative exists for every choice of admissible charts (X, x) and (Y, y), we can say

that the derivative of Df (x) of f at x is defined and given by

(Y,y)

Df (x) = dyy−1 ◦ Dg(X,x) (x(x)) ◦ dxx

(it can be shown that this is well defined and independent of the choice of charts).

**Note that the derivative is now a map between the tangent spaces of the two manifolds
**

Df (x) : Tx V → Ty W . Because of this a common notation for the derivative of f at x is

Tx f . Another alternative notation for the derivative is f∗,x because of its connection to the

category-theoretical pushforward.

Version: 15 Owner: igor Author(s): igor

371.3 l’Hpital’s rule

**L’Hôpital’s rule states that given an unresolvable limit of the form 00 or ∞ ∞
**

, the ratio of

f (x) f 0 (x)

functions g(x) will have the same limit at c as the ratio g0 (x) . In short, if the limit of a ratio

of functions approaches an indeterminate form, then

f (x) f 0 (x)

lim = lim 0

x→c g(x) x→c g (x)

**provided this last limit exists. L’Hôpital’s rule may be applied indefinitely as long0 as the
**

conditions still exist. However it is important to note, that the nonexistance of lim fg0 (x)

(x)

does

not prove the nonexistance of lim fg(x)

(x)

.

Example: We try to determine the value of

x2

lim .

x→∞ ex

1460

∞

As x approaches ∞ the expression becomes an indeterminate form ∞

. By applying L’Hôpital’s

rule we get

x2 2x 2

lim x = lim x = lim x = 0.

x→∞ e x→∞ e x→∞ e

Version: 8 Owner: mathwizard Author(s): mathwizard, slider142

371.4 proof of De l’Hpital’s rule

**Let x0 ∈ R, I be an interval containing x0 and let f and g be two differentiable functions
**

defined on I \ {x0 } with g 0(x) 6= 0 for all x ∈ I. Suppose that

lim f (x) = 0, lim g(x) = 0

x→x0 x→x0

and that

f 0 (x)

lim = m.

x→x0 g 0 (x)

**We want to prove that hence g(x) 6= 0 for all x ∈ I \ {x0 } and
**

f (x)

lim = m.

x→x0 g(x)

**First of all (with little abuse of notation) we suppose that f and g are defined also in the
**

point x0 by f (x0 ) = 0 and g(x0 ) = 0. The resulting functions are continuous in x0 and hence

in the whole interval I.

**Let us first prove that g(x) 6= 0 for all x ∈ I \ {x0 }. If by contradiction g(x̄) = 0 since we
**

also have g(x0 ) = 0, by Rolle’s theorem we get that g 0 (ξ) = 0 for some ξ ∈ (x0 , x̄) which is

against our hypotheses.

**Consider now any sequence xn → x0 with xn ∈ I \ {x0 }. By Cauchy’s mean value theorem
**

there exists a sequence x0n such that

f (xn ) f (xn ) − f (x0 ) f 0 (x0 )

= = 0 0n .

g(xn ) g(xn ) − g(x0 ) g (xn )

But as xn → x0 and since x0n ∈ (x0 , xn ) we get that x0n → x0 and hence

f (xn ) f 0 (xn ) f 0 (x)

lim = lim 0 = lim 0 = m.

n→∞ g(xn ) n→∞ g (xn ) x→x0 g (x)

**Since this is true for any given sequence xn → x0 we conclude that
**

f (x)

lim = m.

x→x0 g(x)

Version: 5 Owner: paolini Author(s): paolini

1461

371.5 related rates

**The notion of a derivative has numerous interpretations and applications. A well-known
**

geometric interpretation is that of a slope, or more generally that of a linear approximation

to a mapping between linear spaces (see here). Another useful interpretation comes from

physics and is based on the idea of related rates. This second point of view is quite general,

and sheds light on the definition of the derivative of a manifold mapping (the latter is

described in the pushforward entry).

Consider two physical quantities x and y that are somehow coupled. For example:

**• the quantities x and y could be the coordinates of a point as it moves along the
**

unit circle;

• the quantity x could be the radius of a sphere and y the sphere’s surface area;

**• the quantity x could be the horizontal position of a point on a given curve and y the
**

distance traversed by that point as it moves from some fixed starting position;

**• the quantity x could be depth of water in a conical tank and y the rate at which the
**

water flows out the bottom.

**Regardless of the application, the situation is such that a change in the value of one quantity
**

is accompanied by a change in the value of the other quantity. So let’s imagine that we

take control of one of the quantities, say x, and change it in any way we like. As we do so,

quantity y follows suit and changes along with x. Now the analytical relation between the

values of x and y could be quite complicated and non-linear, but the relation between the

instantaneous rates of change of x and y is linear.

It does not matter how we vary the two quantities, the ratio of the rates of change depends

only on the values of x and y. This ratio is, of course, the derivative of the function that

maps the values of x to the values of y. Letting ẋ, ẏ denote the rates of change of the two

quantities, we describe this conception of the derivative as

dy ẏ

= ,

dx ẋ

or equivalently as

dy

ẏ = ẋ. (371.5.1)

dx

**Next, let us generalize the discussion and suppose that the two quantities x and y represent
**

physical states with multiple degrees of freedom. For example, x could be a point on the

earth’s surface, and y the position of a point 1 kilometer to the north of x. Again, the

dependence of y and x is, in general, non-linear, but the rate of change of y does have a

linear dependence on the rate of change of x. We would like to say that the derivative is

1462

precisely this linear relation, but we must first contend with the following complication. The

rates of change are no longer scalars, but rather velocity vectors, and therefore the derivative

must be regarded as a linear transformation that changes one vector into another.

**In order to formalize this generalized notion of the derivative we must consider x and y
**

to be points on manifolds X and Y , and the relation between them a manifold mapping

φ : X → Y . A varying x is formally described by a trajectory

γ : I → X, I ⊂ R.

The corresponding velocities take their value in the tangent spaces of X:

γ 0 (t) ∈ Tγ(t) X.

The “coupling” of the two quantities is described by the composition

φ ◦ γ : I → Y.

The derivative of φ at any given x ∈ X is a linear mapping

φ∗ (x) : Tx X → Tφ(x) Y,

**called the pushforward of φ at x, with the property that for every trajectory γ passing
**

through x at time t, we have

(φ ◦ γ)0 (t) = φ∗ (x)γ 0 (t).

The above is the multi-dimensional and coordinate-free generalization of the related rates

relation (370.5.1).

**All of the above has a perfectly rigorous presentation in terms of manifold theory. The
**

approach of the present entry is more informal; our ambition was merely to motivate the

notion of a derivative by describing it as a linear transformation between velocity vectors.

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1463

Chapter 372

26A27 – Nondifferentiability

(nondifferentiable functions, points of

nondifferentiability), discontinuous

derivatives

372.1 Weierstrass function

**The Weierstrass function is a continuous function that is nowhere differentiable, and hence
**

is not an analytic function. The formula for the Weierstrass function is

∞

X

f (x) = bn cos(an πx)

n=1

with a odd, 0 < b < 1, and ab > 1 + 23 π.

**Another example of an everywhere continuous but nowhere differentiable curve is the fractal
**

Koch curve.

[insert plot of Weierstrass function]

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1464

Chapter 373

26A36 – Antidifferentiation

373.1 antiderivative

**The function F (x) is called an antiderivative of a function f (x) if (and only if) the
**

derivative of F is equal to f .

F 0 (x) = f (x)

Note that there are an infinite number of antiderivatives for any function f (x), since any

constant can be added or subtracted from any valid antiderivative to yield another equally

valid antiderivative. To account for this, we express the general antiderivative, or indef-

inite integral, as follows:

intf (x) dx = F (x) + C

where C is an arbitrary constant called the constant of integration. The dx portion means

”with respect to x”, because after all, our functions F and f are functions of x.

Version: 4 Owner: xriso Author(s): xriso

373.2 integration by parts

**When one has an integral of a product of two functions, it is sometimes preferable to simplify
**

the integrand by integrating one of the functions and differentiating the other. This process

is called integrating by parts, and is defined in the following way, where u and v are functions

of x.

intu · v 0 dx = u · v − intv · u0 dx

This process may be repeated indefinitely, and in some cases it may be used to solve for the

original integral algebraically. For definite integrals, the rule appears as

intba u(x) · v 0 (x) dx = (u(b) · v(b) − u(a) · v(a)) − intba v(x) · u0 (x) dx

1465

Proof: Integration by parts is simply the antiderivative of a product rule. Let G(x) =

u(x) · v(x). Then,

G0 (x) = u0 (x)v(x) + u(x)v 0 (x)

Therefore,

G0 (x) − v(x)u0 (x) = u(x)v 0 (x)

We can now integrate both sides with respect to x to get

G(x) − intv(x)u0 (x) dx = intu(x)v 0 (x) dx

**which is just integration by parts rearranged.
**

Example: We integrate the function f (x) = x sin x: Therefore we define u(x) := x and

v 0 (x) = sin x. So integration by parts yields us:

intx sin xdx = −x cos x + int cos xdx = −x cos x + sin x.

Version: 5 Owner: mathwizard Author(s): mathwizard, slider142

373.3 integrations by parts for the Lebesgue integral

**Theorem [1, 2] Suppose f, g are complex valued functions on a bounded interval [a, b]. If f
**

and g are absolutely continuous, then

int[a,b] f 0 g = −int[a,b] f g 0 + f (b)g(b) − f (a)g(a).

where both integrals are Lebesgue integrals.

**Remark Any absolutely continuous function can be differentiated almost everywhere. Thus,
**

in the above, the functions f 0 and g 0 make sense.

**Proof. Since f, g and f g are almost everywhere differentiable with Lebesgue integrable
**

derivatives (see this page), we have

(f g)0 = f 0 g + f g 0

**almost everywhere, and
**

int[a,b] (f g)0 = int[a,b] f 0 g + f g 0

= int[a,b] f 0 g + int[a,b] f g 0 .

The last equality is justified since f 0 g and f g 0 are integrable. For instance, we have

**int[a,b] |f 0 g| ≤ max |g(x)|int[a,b] |f 0|,
**

x∈[a,b]

**which is finite since g is continuous and f 0 is Lebesgue integrable. Now the claim follows
**

from the Fundamental theorem of calculus for the Lebesgue integral. 2

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REFERENCES

1. Jones, F., Lebesgue Integration on Euclidean Spaces, Jones and Barlett Publishers, 1993.

2. Ng, Tze Beng, Integration by Parts, online.

Version: 4 Owner: matte Author(s): matte

1467

Chapter 374

**26A42 – Integrals of Riemann,
**

Stieltjes and Lebesgue type

374.1 Riemann sum

**Suppose there is a function f : I → R where I = [a, b] is a closed interval, and f is bounded
**

on I. If we have a finite set of points {x0 , x1 , x2 , . . . xn } such that a = x0 < x1 < x2 · · · <

xn = b, then this set creates a partition P = {[x0 , x1 ), [x1 , x2 ), . . . [xn − 1, xn ]} of I. If P is a

partition with n ∈ N elements of I, then the Riemann sum of f over I with the partition P

is defined as

n

X

S= f (yi)(xi − xi−1 )

i=1

**where xi−1 6 yi 6 xi . The choice of yi is arbitrary. If yi = xi−1 for all i, then S is called a
**

left Riemann sum. If yi = xi , then S is called a right Riemann sum. Suppose we have

n

X

S= b(xi − xi−1 )

i=1

**where b is the supremum of f over [xi−1 , xi ]; then S is defined to be an upper Riemann sum.
**

Similarly, if b is the infimum of f over [xi−1 , xi ], then S is a lower Riemann sum.

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374.2 Riemann-Stieltjes integral

**Let f and α be bounded, real-valued functions defined upon a closed finite interval I = [a, b]
**

of R(a 6= b), P = {x0 , ..., xn } a partition of I, and ti a point of the subinterval [xi−1 , xi ]. A

sum of the form

n

X

S(P, f, α) = f (ti )(α(xi ) − α(xi−1 ))

i=1

**is called a Riemann-Stieltjes sum of f with respect to α. f is said to be Riemann integrable
**

with respect to α on I if there exists A ∈ R such that given any > 0 there exists a par-

tition P of I for which, for all P finer than P and for every choice of points ti , we have

|S(P, f, α) − A| <

**If such an A exists, then it is unique and is known as the Riemann-Stieltjes integral of
**

f with respect to α. f is known as the integrand and α the integrator. The integral is

denoted by

intba f dα or intba f (x)dα(x)

Version: 3 Owner: vypertd Author(s): vypertd

374.3 continuous functions are Riemann integrable

Let f : [a, b] → R be a continuous function. Then f is Riemann integrable.

Version: 2 Owner: paolini Author(s): paolini

374.4 generalized Riemann integral

**A function f : [a, b] → R is said to be generalized Riemann integrable on [a, b] if there
**

exists a number L ∈ R such that for every > 0 there exists a gauge δ on [a, b] such that if

Ṗ is any δ -fine partition of [a, b], then

|S(f ; Ṗ) − L| <

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Where S(f ; Ṗ) is the Riemann sum for f using the partition Ṗ. The collection of all gener-

alized Riemann integrable functions is usually denoted by R∗ [a, b].

**If f ∈ R∗ [a, b] then the number L is uniquely determined, and is called the generalized
**

Riemann integral of f over [a, b].

Version: 3 Owner: vypertd Author(s): vypertd

**374.5 proof of Continuous functions are Riemann in-
**

tegrable

**Recall the definition of Riemann integral. To prove that f is integrable we have to prove
**

that limδ→0+ S ∗ (δ) − S∗ (δ) = 0. Since S ∗ (δ) is decreasing and S∗ (δ) is increasing it is enough

to show that given > 0 there exists δ > 0 such that S ∗ (δ) − S∗ (δ) < .

So let > 0 be fixed.

**By Heine-Cantor theorem f is uniformly continuous i.e.
**

∃δ > 0 |x − y| < δ ⇒ |f (x) − f (y)| < .

b−a

**Let now P be any partition of [a, b] in C(δ) i.e. a partition {x0 = a, x1 , . . . , xN = b} such
**

that xi+1 − xi < δ. In any small interval [xi , xi+1 ] the function f (being continuous) has a

maximum Mi and minimum mi . Being f uniformly continuous and being xi+1 − xi < δ we

hence have Mi − mi < /(b − a). So the difference between upper and lower Riemann sums

is X X X

Mi (xi+1 − xi ) − mi (xi+1 − xi ) ≤ (xi+1 − xi ) = .

i i

b − a i

Being this true for every partition P in C(δ) we conclude that S ∗ (δ) − S∗ (δ) < .

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1470

Chapter 375

26A51 – Convexity, generalizations

375.1 concave function

**Let f (x) a continuous function defined on an interval [a, b]. Then we say that f is a concave
**

function on [a, b] if, for any x1 , x2 in [a, b] and any λ ∈ [0, 1] we have

f λx1 + (1 − λ)x2 > λf (x1 ) + (1 − λ)f (x2 ).

The definition is equivalent to the statements:

• For all x1 , x2 in [a, b],

x1 + x2 f (x1 ) + f (x2 )

f ≥

2 2

• The second derivative of f is negative on [a, b].

• If f has a derivative which is monotone decreasing.

obviously, the last two items apply provided f has the required derivatives.

And example of concave function is f (x) = −x2 on the interval [−5, 5].

Version: 5 Owner: drini Author(s): drini

1471

Chapter 376

26Axx – Functions of one variable

376.1 function centroid

**Let f : D ⊂ R → R be an arbitrary function. By analogy with the geometric centroid, the
**

centroid of an function f is defined as:

intxf (x)dx

hxi = ,

intf (x)dx

where the integrals are taken over the domain D.

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1472

Chapter 377

**26B05 – Continuity and
**

differentiation questions

377.1 C0∞ (U ) is not empty

**Theorem If U is a non-empty open set in Rn , then the set of smooth functions with compact support
**

C0∞ (U) is not empty.

The proof is divided into three sub-claims:

**Claim 1 Let a < b be real numbers. Then there exists a smooth non-negative function
**

f : R → R, whose support is the compact set [a, b].

To prove Claim 1, we need the following lemma:

Lemma ([4], pp. 14) If

0 for x ≤ 0,

φ(x) =

e−1/x for x > 0,

then φ : R → R is a non-negative smooth function.

(A proof of the Lemma can be found in [4].)

Proof of Claim 1. Using the lemma, let us define

f (x) = φ(x − a)φ(b − x).

**Since φ is smooth, it follows that f is smooth. Also, from the definition of φ, we see that
**

φ(x − a) = 0 precisely when x ≤ a, and φ(b − x) = 0 precisely when x ≥ b. Thus the support

of f is indeed [a, b]. 2

Claim 2 Let ai , bi be real numbers with ai < bi for all i = 1, . . . , n. Then there exists a

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smooth non-negative function f : Rn → R whose support is the compact set [a1 , b1 ] × · · · ×

[an , bn ].

**Proof of Claim 2. Using Claim 1, we can for each i = 1, . . . , n construct a function fi with
**

support in [ai , bi ]. Then

f (x1 , . . . , xn ) = f1 (x1 )f2 (x2 ) · · · fn (xn )

gives a smooth function with the sought properties. 2

**Claim 3 If U is a non-empty open set in Rn , then there are real numbers ai < bi for
**

i = 1, . . . , n such that [a1 , b1 ] × · · · × [an , bn ] is a subset of U.

**Proof of Claim 3. Here, of course, we assume that Rn is equipped with the usual topology
**

induced by the open balls of the Euclidean metric.

**Since U is non-empty, there exists some point x in U. Further, since U is a topological space,
**

x is contained in some open set. Since the topology has a basis consisting of open balls, there

exists a y ∈ U and ε > 0 such that x is contained in the open ball B(y, ε). Let us now set

ai = yi − 2√ε n and bi = yi + 2√ε n for all i = 1, . . . , n. Then D = [a1 , b1 ] × · · · × [an , bn ] can be

parametrized as

ε

D = {y + (λ1 , . . . , λn ) √ | λi ∈ [−1, 1] for all i = 1, . . . , n}.

2 n

For an arbitrary point in D, we have

ε ε

|y + (λ1 , . . . , λn ) √ − y| = |(λ1 , . . . , λn )| √ |

2 n 2 n

q

ε

= √ λ21 + · · · + λ2n

2 n

ε

≤ < ε,

2

so D ⊂ B(y, ) ⊂ U, and Claim 3 follows. 2

REFERENCES

1. L. Hörmander, The Analysis of Linear Partial Differential Operators I, (Distribution

theory and Fourier Analysis), 2nd ed, Springer-Verlag, 1990.

Version: 3 Owner: matte Author(s): matte

377.2 Rademacher’s Theorem

**Let f : Rn → R be any Lipschitz continuous function. Then f is differentiable in almost
**

every x ∈ Rn .

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Version: 1 Owner: paolini Author(s): paolini

377.3 smooth functions with compact support

**Definition [3] Let U be an open set in Rn . Then the set of smooth functions with
**

compact support (in U) is the set of functions f : Rn → C which are smooth (i.e.,

∂ α f : Rn → C is a continuous function for all multi-indices α) and supp f is compact and

contained in U. This functionspace is denoted by C0∞ (U).

Remarks

**1. A proof that C0∞ (U) is not empty can be found here.
**

2. With the usual point-wise addition and point-wise multiplication by a scalar, C0∞ (U)

is a vector space over the field C.

3. Suppose U and V are open subsets in Rn and U ⊂ V . Then C0∞ (U) is a vector subspace

of C0∞ (V ). In particular, C0∞ (U) ⊂ C0∞ (V ).

It is possible to equip C0∞ (U) with a topology, which makes C0∞ (U) into a locally convex topological vector s

The definition, however, of this topology is rather involved (see e.g. [3]). However, the next

theorem shows when a sequence converges in this topology.

**Theorem 1 Suppose that U is an open set in Rn , and that {φi }∞
**

i=1 is a sequence of functions

∞

in C0 (U). Then {φi } converges (in the aforementioned topology) to a function φ ∈ C0∞ (U)

if and only if the following conditions hold:

**1. There is a compact set K ⊂ U such that supp φi ⊂ K for all i = 1, 2, . . ..
**

2. For every multi-index α,

∂ α φi → ∂ α φ

in the sup-norm.

**Theorem 2 Suppose that U is an open set in Rn , that Γ is a locally convex topological
**

vector space, and that L : C0∞ (U) → Γ is a linear map. Then L is a continuous map, if and

only if the following condition holds:

**If K is a compact subset of U, and {φi }∞ ∞
**

i=1 is a sequence of functions in C0 (U) such

that supp φi ⊂ K for all i, and φi → φ (in C0∞ (U)) for some φ ∈ D(U), then Lφi → Lφ

(in C).

The above theorems are stated without proof in [1].

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REFERENCES

1. W. Rudin, Functional Analysis, McGraw-Hill Book Company, 1973.

2. G.B. Folland, Real Analysis: Modern Techniques and Their Applications, 2nd ed,

John Wiley & Sons, Inc., 1999.

Version: 3 Owner: matte Author(s): matte

1476

Chapter 378

**26B10 – Implicit function theorems,
**

Jacobians, transformations with

several variables

378.1 Jacobian matrix

**The Jacobian [Jf~(x)] of a function f~ : Rn → Rm is the matrix of partial derivatives such
**

that

D1 f1 (x) . . . Dn f1 (x)

.. .. ..

[Jf~(x)] = . . .

D1 fm (x) . . . Dn fm (x)

A more concise way of writing it is

~ 1

∇f

−−→ −−→

~

[Jf(x)] = [D1 f , · · · , Dn f ] = ...

~ m

∇f

−−→ ~ m is the gradient

where Dn f is the partial derivative with respect to the nth variable and ∇f

of the nth component of f. The Jacobian matrix represents the full derivative matrix [Df (x)]

of f at x iff f is differentiable at x. Also, if f is differentiable at x, then [Jf(x)] = [Df (x)]

and the directional derivative in the direction ~v is [Df(x)]~v .

Version: 9 Owner: slider142 Author(s): slider142

378.2 directional derivative

Partial derivatives measure the rate at which a multivariable function f~ varies as the variable

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moves in the direction of the standard basis vectors. Directional derivatives measure the rate

at which f~ varies when the variable moves in the direction ~v. Thus the directional derivative

of f~ at a in the direction ~v is represented as

∂ f~(a) f~(a + h~v ) − f~(a)

D~v f (a) = = lim .

∂~v h→0 h

x

For example, if f y = x2 + 3y 2z, and we wanted to find the derivative at the point

z

1 1

a = 2 in the direction ~v = 1 , our equation would be

3 1

limh→0 h1 ((1 + h)2 + 3(2 + h)2 (3 + h) − 37)

= limh→0 h1 (3h3 + 37h2 + 50h) .

**= limh→0 3h2 + 37h + 50 = 50
**

One may also use the Jacobian matrix if the function is differentiable to find the derivative

in the direction ~v as [Jf (x)]~v .

Version: 6 Owner: slider142 Author(s): slider142

378.3 gradient

**Summary. The gradient is a first-order differential operator that maps functions to vector fields.
**

It is a generalization of the ordinary derivative, and as such conveys information about the

rate of change of a function relative to small variations in the independent variables. The

gradient of a function f is customarily denoted by ∇f or by grad f .

**Definition: Euclidean space Consider n-dimensional Euclidean space with orthogonal
**

coordinates x1 , . . . , xn , and corresponding unit vectors e1 , . . . , en . In this setting, the gradi-

ent of a function f (x1 , . . . , xn ) is defined to be the vector field given by

Xn

∂f

∇f = ei .

i=1

∂xi

**It is also useful to represent the gradient operator as the vector-valued differential operator
**

n

X ∂

∇= ei .

i=1

∂xi

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or, in the context of Euclidean 3-space, as

∂ ∂ ∂

∇=i +j +k ,

∂x ∂y ∂z

where i, j, k are the unit vectors lying along the positive direction of the x, y, z axes, respec-

tively. Using this formalism, the ∇ symbol can be used to express the divergence operator

as ∇·, the curl operator as ∇×, and the Laplacian operator as ∇2 . To wit, for a given vector

field

A = Ax i + Ay j + Az k,

and a given function f we have

∂Ax ∂Ay ∂Az

∇·A= + +

∂x ∂y ∂z

∂Az ∂Ay ∂Ax ∂Az ∂Ay ∂Ax

∇×A= − i+ − j+ − k

∂y ∂z ∂z ∂x ∂x ∂y

∂2f ∂2f ∂2f

∇2 f = + + .

∂x2 ∂y 2 ∂z 2

**Definition: Riemannian geometry More generally still, consider a Riemannian manifold
**

with metric tensor gij and inverse g ij . In this setting the gradient X = grad f of a function

f relative to a general coordinate system, is given by

X j = g ij f,i . (378.3.1)

Note that the Einstein summation convention is in force above. Also note that f,i denotes

the partial derivative of f with respect to the ith coordinate.

**Definition (377.3.1) is useful even in the Euclidean setting, because it can be used to derive
**

the formula for the gradient in various generalized coordinate systems. For example, in the

cylindrical system of coordinates (r, θ, z) we have

1 0 0

gij = 0 r 2 0

0 0 1

while for the system of spherical coordinates (ρ, φ, θ) we have

1 0 0

gij = 0 ρ2 0 .

0 0 ρ sin2 φ

2

**Hence, for a given function f we have
**

∂f 1 ∂f ∂f

∇f = er + eθ + k Cylindrical

∂r r ∂θ ∂z

∂f 1 ∂f 1 ∂f

∇f = eρ + eφ + eθ Spherical ,

∂ρ ρ ∂φ ρ sin φ ∂θ

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where for the cylindrical system

∂ x y

er = = i+ j

∂r r r

1 ∂ y x

eθ = =− i+ j

r ∂θ r r

are the unit vectors in the direction of increase of r and θ, respectively, and for the spherical

system

∂ x y z

eρ = = i+ j+ k

∂ρ ρ ρ ρ

1 ∂ zx zy r

eφ = = i+ j− k

ρ ∂φ rρ rρ ρ

1 ∂ y x

eθ = =− i+ j

ρ sin θ ∂θ r r

are the unit vectors in the direction of increase of ρ, θ, φ, respectively.

**Physical Interpretation. In the simplest case, we consider the Euclidean plane with
**

Cartesian coordinates x, y. The gradient of a function f (x, y) is given by

∂f ∂f

∇f = i+ j,

∂x ∂y

where i, j denote, respectively, the standard unit horizontal and vertical vectors. The gradient

vectors have the following geometric interpretation. Consider the graph z = f (x, y) as a

surface in 3-space. The direction of the gradient vector ∇f is the direction of steepest

ascent, while the magnitude is the slope in that direction. Thus,

s 2

2

∂f ∂f

k∇f k = +

∂x ∂y

describes the steepness of the hill z = f (x, y) at a point on the hill located at (x, y, f (x, y)).

**A more general conception of the gradient is based on the interpretation of a function f as
**

a potential corresponding to some conservative physical force. The negation of the gradient,

−∇f , is then interpreted as the corresponding force field.

**Differential identities. Several properties of the one-dimensional derivative generalize to
**

a multi-dimensional setting

∇(af + bg) = a∇f + b∇g Linearity

∇(f g) = f ∇g + g∇f Product rule

∇(φ ◦ f ) = (φ0 ◦ f )∇f Chain rule

Version: 9 Owner: rmilson Author(s): rmilson, slider142

1480

378.4 implicit differentiation

**Implicit differentiation is a tool used to analyze functions that cannot be conveniently put
**

into a form y = f (x) where x = (x1 , x2 , ..., xn ). To use implicit differentiation meaningfully,

you must be certain that your function is of the form f (x) = 0 (it can be written as

a level set) and that it satisfies the implicit function theorem (f must be continuous, its

first partial derivatives must be continuous, and the derivative with respect to the implicit

function must be non-zero). To actually differentiate implicitly, we use the chain rule to

differentiate the entire equation.

**Example: The first step is to identify the implicit function. For simplicity in the example,
**

we will assume f (x, y) = 0 and y is an implicit function of x. Let f (x, y) = x2 + y 2 + xy = 0

(Since this is a two dimensional equation, all one has to check is that the graph of y may

be an implicit function of x in local neighborhoods.) Then, to differentiate implicitly, we

differentiate both sides of the equation with respect to x. We will get

dy dy

2x + 2y · +x·1· +y =0

dx dx

Do you see how we used the chain rule in the above equation ? Next, we simply solve for our

dy

implicit derivative dx = − 2x+y

2y+x

. Note that the derivative depends on both the variable and

the implicit function y. Most of your derivatives will be functions of one or all the variables,

including the implicit function itself.

[better example and ?multidimensional? coming]

Version: 2 Owner: slider142 Author(s): slider142

378.5 implicit function theorem

**Let f = (f1 , ..., fn ) be a continuously differentiable, vector-valued function mapping an
**

open set E ⊂ Rn+m into Rn . Let (a, b) = (a1 , ..., an , b1 , ..., bm ) be a point in E for which

f(a, b) = 0 and such that the n × n determinant

|Dj fi (a, b)| =

6 0

**for i, j = 1, ..., n. Then there exists an m-dimensional neighbourhood W of b and a unique
**

continuously differentiable function g : W → Rn such that g(b) = a and

f(g(t), t) = 0

for all t ∈ W .

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Simplest case

**When n = m = 1, the theorem reduces to: Let F be a continuously differentiable, real-
**

valued function defined on an open set E ⊂ R2 and let (x0 , y0) be a point on E for which

F (x0 , y0) = 0 and such that

∂F

|x ,y 6= 0

∂x 0 0

Then there exists an open interval I containing y0 , and a unique function f : I → R which

is continuously differentiable and such that f (y0 ) = x0 and

F (f (y), y) = 0

for all y ∈ I.

Note

**The inverse function theorem is a special case of the implicit function theorem where the
**

dimension of each variable is the same.

Version: 7 Owner: vypertd Author(s): vypertd

378.6 proof of implicit function theorem

Consider the function F : E → Rn × Rm defined by

F (x, y) = (f (x, y), y).

∂f j ∂fj

Setting Ajk = ∂xk

(a, b), and Mji = ∂yi

(a, b), A is an n × m matrix and M is n × n. It holds

Df (a, b) = (A|M)

and hence

In 0

DF (a, b) = .

A M

**Being det M 6= 0 M is invertible and hence DF (a, b) is invertible too. Applying the
**

inverse function theorem to F we find that there exist a neighbourhood V of a and W of b

and a function G ∈ C 1 (V × W, Rn+m ) such that F (G(x, y)) = (x, y) for all (x, y) ∈ V × W .

Letting G(x, y) = (G1 (x, y), G2(x, y)) (so that G1 : V × W → Rn , G2 : V × W → Rm ) we

hence have

(x, y) = F (G1 (x, y), G2(x, y)) = (f (G1 (x, y), G2(x, y)), G2 (x, y))

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and hence y = G2 (x, y) and x = f (G1 (x, y), G2(x, y)) = f (G1 (x, y), y). So we only have to

set g(y) = G1 (0, y) to obtain

f (g(y), y) = 0, ∀y ∈ W.

Version: 1 Owner: paolini Author(s): paolini

1483

Chapter 379

26B12 – Calculus of vector functions

379.1 Clairaut’s theorem

**Theorem. (Clairaut’s Theorem) If F : Rn → Rm is a function whose second partial derivatives
**

exist and are continuous on a set S ⊆ Rn , then

∂2f ∂2f

=

∂xi ∂xj ∂xj ∂xi

on S (where 1 6 i, j 6 n).

**This theorem is commonly referred to as simply ’the equality of mixed partials’. It is usually
**

first presented in a vector calculus course, and is useful in this context for proving basic

properties of the interrelations of gradient, divergence, and curl. I.e., if F : R3 → R3 is a

function satisfying the hypothesis, then ∇ · (∇ × F) = 0. Or, if f : R3 → R is a function

satisfying the hypothesis, ∇ × ∇f = 0.

Version: 10 Owner: flynnheiss Author(s): flynnheiss

379.2 Fubini’s Theorem

**Fubini’s Theorem Let I ⊂ RN and J ⊂ RM be compact intervals, and let f : I × J → RK
**

be a Riemann integrable function such that, for each x ∈ I the integral

F (x) := intJ f (x, y) dµJ (y)

exists. Then F : I → RK is Riemann integrable, and

intI F = intI×J f.

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This theorem effectively states that, given a function of N variables, you may integrate it

one variable at a time, and that the order of integration does not affect the result.

**Example Let I := [0, π/2] × [0, π/2], and let f : I → R, x 7→ sin(x) cos(y) be a function.
**

Then

ZZ

intI f = sin(x) cos(y)

[0,π/2]×[0,π/2]

π/2 π/2

= int0 int0 sin(x) cos(y) dy dx

π/2

= int0 sin(x) (1 − 0) dx = (0 − −1) = 1

R R

Note that it is often simpler (and no less correct) to write ··· I

f as intI f .

Version: 3 Owner: vernondalhart Author(s): vernondalhart

379.3 Generalised N-dimensional Riemann Sum

**Let I = [a1 , b1 ] × · · · × [aN , bN ] be an N-cell in RN . For each j = 1, . . . , N, let aj = tj,0 <
**

. . . < tj,N = bj be a partition Pj of [aj , bj ]. We define a partition P of I as

P := P1 × · · · × PN

Each partition P of I generates a subdivision of I (denoted by (Iν )ν ) of the form

Iν = [t1,j , t1,j+1 ] × · · · × [tN,k , tN,k+1 ]

**Let f : U → RM be such that I ⊂ U, and let (Iν )ν be the corresponding subdivision of a
**

partition P of I. For each ν, choose xν ∈ Iν . Define

X

S(f, P ) := f (xν )µ(Iν)

ν

As the Riemann sum of f corresponding to the partition P .

A partition Q of I is called a refinement of P if P ⊂ Q.

Version: 1 Owner: vernondalhart Author(s): vernondalhart

379.4 Generalized N-dimensional Riemann Integral

**Let I = [a1 , b1 ] × · · · × [aN , bN ] ⊂ RN be a compact interval, and let f : I → RM be a
**

function. Let > 0. If there exists a y ∈ RM and a partition P of I such that for each

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refinement P of P (and corresponding Riemann sum S(f, P )),

kS(f, P ) − yk <

Then we say that f is Riemann integrable over I, that y is the Riemann integral of f over

I, and we write

intI f := intI f dµ := y

**Note also that it is possible to extend this definition to more arbitrary sets; for any bounded
**

set D, one can find a compact interval I such that D ⊂ I, and define a function

(

f (x), x ∈ D

f˜ : I → RM x 7→

0, x∈/D

in which case we define

intD f := intI f˜

Version: 3 Owner: vernondalhart Author(s): vernondalhart

379.5 Helmholtz equation

**It is a partial differential equation which, in scalar form is
**

∇2 f + k 2 f = 0,

or in vector form is

∇2 A + k 2 A = 0,

where ∇2 is the Laplacian. The solutions of this equation represent the solution of the

wave equation, which is of great interest in physics.

**Consider a wave equation
**

∂2ψ

= c 2 ∇2 ψ

∂t2

with wave speed c. If we look for time harmonic standing waves of frequency ω,

ψ(x, t) = e−jωt φ(x)

we find that φ(x) satisfies the Helmholtz equation:

(∇2 + k 2 )φ = 0

where k = ω/c is the wave number.

**Usually Helmholtz equation is solved by seperation of variables method, in cartesian, spher-
**

ical or cylindrical coordinates.

Version: 3 Owner: giri Author(s): giri

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379.6 Hessian matrix

**The Hessian of a scalar function of a vector is the matrix of partial second derivatives. So
**

the Hessian matrix of a function f : Rn → R is:

** ∂2f ∂2f ∂2f
**

dx21 dx1 dx2

... dx1 dxn

∂2f ∂2f ∂2f

dx22

...

dx2 dx1 dx2 dxn

(379.6.1)

.. .. .. ..

. . . .

∂2f ∂2f ∂2f

dxn dx1 dxn dx2

... dx2n

Note that the Hessian is symmetric because of the equality of mixed partials.

Version: 2 Owner: bshanks Author(s): akrowne, bshanks

379.7 Jordan Content of an N-cell

**Let I = [a1 , b1 ] × · · · × [aN , bN ] be an N-cell in RN . Then the Jordan content (denoted µ(I))
**

of I is defined as

YN

µ(I) := (bj − aj )

j=1

Version: 1 Owner: vernondalhart Author(s): vernondalhart

379.8 Laplace equation

The scalar form of Laplace’s equation is the partial differential equation

∇2 f = 0

**and the vector form is
**

∇2 A = 0,

where ∇2 is the Laplacian. It is a special case of the Helmholtz differential equation with

k = 0.

**A function f which satisfies Laplace’s equation is said to be harmonic. Since Laplace’s
**

equation is linear, the superposition of any two solutions is also a solution.

Version: 3 Owner: giri Author(s): giri

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379.9 chain rule (several variables)

**The chain rule is a theorem of analysis that governs derivatives of composed functions. The
**

basic theorem is the chain rule for functions of one variables (see here). This entry is devoted

to the more general version involving functions of several variables and partial derivatives.

Note: the symbol Dk will be used to denote the partial derivative with respect to the k th

variable.

**Let F (x1 , . . . , xn ) and G1 (x1 , . . . , xm ), . . . , Gn (x1 , . . . , xm ) be differentiable functions of sev-
**

eral variables, and let

H(x1 , . . . , xm ) = F (G1 (x1 , . . . , xm ), . . . , Gn (x1 , . . . , xm ))

**be the function determined by the composition of F with G1 , . . . , Gn The partial derivatives
**

of H are given by

n

X

(Dk H)(x1 , . . . , xm ) = (Di F )(G1 (x1 , . . . , xm ), . . .)(Dk Gi )(x1 , . . . , xm ).

i=1

**The chain rule can be more compactly (albeit less precisely) expressed in terms of the Jacobi-
**

Legendre partial derivative symbols (historical note). Just as in the Leibniz system, the basic

idea is that of one quantity (i.e. variable) depending on one or more other quantities. Thus

we would speak about a variable z depends differentiably on y1 , . . . , yn , which in turn depend

differentiably on variables x1 , . . . , xm . We would then write the chain rule as

n

X ∂z ∂yi

∂z

= , j = 1, . . . m.

∂xj i=1

∂y i ∂xj

**The most general, and conceptually clear approach to the multi-variable chain is based on the
**

notion of a differentiable mapping, with the Jacobian matrix of partial derivatives playing

the role of generalized derivative. Let, X ⊂ Rm and Y ⊂ Rn be open domains and let

F : Y → Rl , G:X→Y

**be differentiable mappings. In essence, the symbol F represents l functions of n variables
**

each:

F = (F1 , . . . , Fl ), Fi = Fi (x1 , . . . , xn ),

whereas G = (G1 , . . . , Gn ) represents n functions of m variables each. The derivative of such

mappings is no longer a function, but rather a matrix of partial derivatives, customarily called

the Jacobian matrix. Thus

D1 F1 . . . Dn F1 D1 G1 . . . Dm G1

.. ..

DF = ... ..

. . DG = ... ..

. .

D1 Fl . . . Dn Fl D1 Gn . . . Dm Gn

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The chain rule now takes the same form as it did for functions of one variable:

D(F ◦ G) = ((DF) ◦ G) (DG),

albeit with matrix multiplication taking the place of ordinary multiplication.

This form of the chain rule also generalizes quite nicely to the even more general setting

where one is interested in describing the derivative of a composition of mappings between

manifolds.

Version: 7 Owner: rmilson Author(s): rmilson

379.10 divergence

**Basic Definition. Let x, y, z be a system of Cartesian coordinates on 3-dimensional Euclidean space,
**

and let i, j, k be the corresponding basis of unit vectors. The divergence of a continuously differentiable

vector field

F = F 1 i + F 2 j + F 3 k,

is defined to be the function

∂F 1 ∂F 2 ∂F 3

div F = + + .

∂x ∂y ∂z

Another common notation for the divergence is ∇ · F (see gradient), a convenient mnemonic.

**Physical interpretation. In physical terms, the divergence of a vector field is the extent
**

to which the vector field flow behaves like a source or a sink at a given point. Indeed, an

alternative, but logically equivalent definition, gives the divergence as the derivative of the

net flow of the vector field across the surface of a small sphere relative to the surface area of

the sphere. To wit,

(div F)(p) = lim intS(F · N)dS / 4πr 2 ,

r→0

where S denotes the sphere of radius r about a point p ∈ R3 , and the integral is a surface

integral taken with respect to N, the normal to that sphere.

**The non-infinitesimal interpretation of divergence is given by Gauss’s Theorem. This the-
**

orem is a conservation law, stating that the volume total of all sinks and sources, i.e. the

volume integral of the divergence, is equal to the net flow across the volume’s boundary. In

symbols,

intV div F dV = intS (F · N) dS,

where V ⊂ R3 is a compact region with a smooth boundary, and S = ∂V is that boundary

oriented by outward-pointing normals. We note that Gauss’s theorem follows from the more

general Stokes’ theorem, which itself generalizes the fundamental theorem of calculus.

**In light of the physical interpretation, a vector field with constant zero divergence is called
**

incompressible – in this case, no net flow can occur across any closed surface.

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General definition. The notion of divergence has meaning in the more general setting of

Riemannian geometry. To that end, let V be a vector field on a Riemannian manifold. The

covariant derivative of V is a type (1, 1) tensor field. We define the divergence of V to be the

trace of that field. In terms of coordinates (see tensor and Einstein summation convention),

we have

div V = V i ;i .

Version: 6 Owner: rmilson Author(s): rmilson, jaswenso

379.11 extremum

**Extrema are minima and maxima. The singular forms of these words are extremum, mini-
**

mum, and maximum.

**Extrema may be “global” or “local”. A global minimum of a function f is the lowest value
**

that f ever achieves. If you imagine the function as a surface, then a global minimum is the

lowest point on that surface. Formally, it is said that f : U → V has a global minimum at x

if ∀u ∈ U, f (x) 6 f (u).

**A local minimum of a function f is a point x which has less value than all points ”next
**

to” it. If you imagine the function as a surface, then a local minimum is the bottom of a

“valley” or “bowl” in the surface somewhere. Formally, it is said that f : U → V has a local

minimum at x if ∃ a neighborhood N of x such that ∀y ∈ N, f (x) 6 f (y).

If you flip the 6 signs above to >, you get the definitions of global and local maxima.

**A ”strict local minima” or ”strict local maxima” means that nearby points are strictly less
**

than or strictly greater than the critical point, rather than 6 or >. For instance, a strict

local minima at x has a neighborhood N such that ∀y ∈ N, (f (x) < f (y) or y = x).

Related concepts are plateau and saddle point.

Finding minima or maxima is an important task which is part of the field of optimization.

Version: 9 Owner: bshanks Author(s): bshanks, bbukh

379.12 irrotational field

**Suppose Ω is an open set in R3 , and V is a vector field with differentiable real (or possibly
**

complex) valued component functions. If ∇ × V = 0, then V is called an irrotional vector

field, or curl free field.

If U and V are irrotational, then U × V is solenoidal.

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Version: 6 Owner: matte Author(s): matte, giri

379.13 partial derivative

**The partial derivative of a multivariable function f is simply its derivative with respect to
**

only one variable, keeping all other variables constant (which are not functions of the variable

in question). The formal definition is

a1

..

.

∂f 1

f (a + h~ei ) − f (a)

Di f (a) = = lim f ai + h − f (a) = lim

∂ai h→0 h .. h→0 h

.

an

where ~ei is the standard basis vector of the ith variable. Since this only affects the ith vari-

able, one can derive the function using common rules and tables, treating all other variables

(which are not functions of ai ) as constants. For example, if f (x) = x2 + 2xy + y 2 + y 3z,

then

∂f

(1) ∂x

= 2x + 2y

∂f

(2) ∂y

= 2x + 2y + 3y 2z

∂f

(3) ∂z

= y3

Note that in equation (1), we treated y as a constant, since we were differentiating with re-

spect to x. d(c∗x)

dx

= c The partial derivative of a vector-valued function f~(x) with respect

−→ ∂ f~

to variable ai is a vector Di f = ∂a i

.

Multiple Partials:

Multiple partial derivatives can be treated just like multiple derivatives. There is an addi-

tional degree of freedom though, as you can compound derivatives with respect to different

variables. For example, using the above function,

∂2f ∂

(4) ∂x2

= ∂x

(2x + 2y) =2

∂2f ∂

(5) ∂z∂y

= ∂z

(2x + 2y + (5)3y 2z) = 3y 2

∂2f ∂

(6) ∂y∂z

= ∂y

(y 3 ) = 3y 2

**D12 is another way of writing ∂x1∂∂x2 . If f (x) is continuous in the neighborhood of x, it can
**

be shown that Dij f (x) = Dji f (x) where i, j are the ith and jth variables. In fact, as long

as an equal number of partials are taken with respect to each variable, changing the order

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of differentiation will produce the same results in the above condition.

Another form of notation is f (a,b,c,...)(x) where a is the partial derivative with respect to the

first variable a times, b is the partial with respect to the second variable b times, etc.

Version: 17 Owner: slider142 Author(s): slider142

379.14 plateau

A plateau of a function is a region where a function has constant value.

**More formally, let U and V be topological spaces. A plateau for a scalar function f : U → V
**

is a path-connected set of points P ⊆ U such that for some y we have

∀p ∈ P, f (p) = y (379.14.1)

**Please take note that this entry is not authoritative. If you know of a more standard definition
**

of ”plateau”, please contribute it, thank you.

Version: 4 Owner: bshanks Author(s): bshanks

379.15 proof of Green’s theorem

**Consider the region R bounded by the closed curve P in a well-connected space. P can be
**

given by a vector valued function F~ (x, y) = (f (x, y), g(x, y)). The region R can then be

described by

∂g ∂f ∂g ∂f

intintR − dA = intintR dA − intintR dA

∂x ∂y ∂x ∂y

The double integrals above can be evaluated separately. Let’s look at

∂g B(y) ∂g

intintR dA = intba intA(y) dxdy

∂x ∂x

Evaluating the above double integral, we get

intba (g(A(y), y) − g(B(y), y)) dy = intba g(A(y), y) dy − intba g(B(y), y) dy

**According to the fundamental theorem of line integrals, the above equation is actually
**

equivalent to the evaluation of the line integral of the function F~1 (x, y) = (0, g(x, y)) over a

path P = P1 + P2 , where P1 = (A(y), y) and P2 = (B(y), y).

I

inta g(A(y), y) dy − inta g(B(y), y) dy = intP1 F~1 · d~t + intP2 F~1 · d~t =

b b

F~1 · d~t

P

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Thus we have I

∂g

intintR dA = F~1 · d~t

∂x P

By a similar argument, we can show that

I

∂f

intintR dA = − F~2 · d~t

∂y P

**where F~2 = (f (x, y), 0). Putting all of the above together, we can see that
**

I I I I

∂g ∂f

intintR − dA = F~1 · d~t + F~2 · d~t = (F~1 + F~2 ) · d~t = (f (x, y), g(x, y)) · d~t

∂x ∂y P P P P

which is Green’s theorem.

Version: 7 Owner: slider142 Author(s): slider142

**379.16 relations between Hessian matrix and local ex-
**

trema

**Let x be a vector, and let H(x) be the Hessian for f at a point x. Let the neighborhood of x
**

be in the domain for f , and let f have continuous partial derivatives of first and second order.

Let ∇f = ~0.

If H(x) is positive definite, then x is a strict local minimum for f .

If x is a local minimum for x, then H(x) is positive semidefinite.

If H(x) is negative definite, then x is a strict local maximum for f .

If x is a local maximum for x, then H(x) is negative semidefinite.

If H(x) is indefinite, x is a nondegenerate saddle point.

**If the case when the dimension of x is 1 (i.e. f : R → R), this reduces to the Second
**

Derivative Test, which is as follows:

**Let the neighborhood of x be in the domain for f , and let f have continuous partial deriva-
**

tives of first and second order. Let f 0 (x) = 0. If f 00 (x) > 0, then x is a strict local minimum.

If f 00 (x) < 0, then x is a strict local maximum.

Version: 6 Owner: bshanks Author(s): bshanks

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379.17 solenoidal field

A solenoidal vector field is one that satisfies

∇·B=0

at every point where the vector field B is defined. Here ∇ · B is the divergence.

**This condition actually implies that there exists a vector A, known as the vector potential,
**

such that

B = ∇ × A.

For a function f satisfying Laplace’s equation

∇2 f = 0,

it follows that ∇f is solenoidal.

Version: 4 Owner: giri Author(s): giri

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Chapter 380

**26B15 – Integration: length, area,
**

volume

380.1 arc length

**Arc length is the length of a section of a differentiable curve. Finding arc length is useful in
**

many applications, for the length of a curve can be attributed to distance traveled, work, etc.

It is commonly represented as S or the differential ds if one is differentiating or integrating

with respect to change in arclength.

**If one knows the vector function or parametric equations of a curve, finding the arc length
**

is simple, as it can be given by the sum of the lengths of the tangent vectors to the curve or

intba |F~ 0(t)| dt = S

**Note that t is an independent parameter. In Cartesian coordinates, arclength can be calcu-
**

lated by the formula p

S = intba 1 + (f 0 (x))2 dx

This formula is derived by viewing arclength as the Riemman sum

n

X p

lim 1 + f 0 (xi ) ∆x

∆x→∞

i=1

**The term being summed is the length of an approximating secant to the curve over the dis-
**

tance ∆x. As ∆x vanishes, the sum approaches the arclength, thus the algorithm. Arclength

can also be derived for polar coordinates from the general formula for vector functions given

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above. The result is p

L = intba r(θ)2 + (r 0 (θ))2 dθ

Version: 5 Owner: slider142 Author(s): slider142

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Chapter 381

**26B20 – Integral formulas (Stokes,
**

Gauss, Green, etc.)

381.1 Green’s theorem

**Green’s theorem provides a connection between path integrals over a well-connected region
**

in the plane and the area of the region bounded in the plane. Given a closed path P bounding

a region R with area A, and a vector-valued function F~ = (f (x, y), g(x, y)) over the plane,

I

F~ · d~x = int

intR [g1 (x, y) − f2 (x, y)]dA

P

where an is the derivative of a with respect to the nth variable.

**Corollary: The closed path integral over a gradient of a function with continuous partial derivatives
**

is always zero. Thus, gradients are conservative vector fields. The smooth function is called

the potential of the vector field.

Proof: The corollary states that

I

~ h · d~x = 0

∇

P

We can easily prove this using Green’s theorem.

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I

~ h · d~x = intintR [g1 (x, y) − f2 (x, y)]dA

∇

P

But since this is a gradient...

intintR [g1 (x, y) − f2 (x, y)]dA = int

intR [h21 (x, y) − h12 (x, y)]dA

Since h12 = h21 for any function with continuous partials, the corollary is proven.

Version: 4 Owner: slider142 Author(s): slider142

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Chapter 382

26B25 – Convexity, generalizations

382.1 convex function

**Definition Suppose Ω is a convex set in a vector space over R (or C), and suppose f is a
**

function f : Ω → R. If for any x, y ∈ Ω and any λ ∈ (0, 1), we have

f λa + (1 − λ)b 6 λf (a) + (1 − λ)f (b),

**we say that f is a convex function. If for any x, y ∈ Ω and any λ ∈ (0, 1), we have
**

f λa + (1 − λ)b > λf (a) + (1 − λ)f (b),

**we say that f is a concave function. If either of the inequalities are strict, then we say
**

that f is a strictly convex function, or a strictly concave function, respectively.

Properties

**• A function f is a (strictly) convex function if and only if −f is a (strictly) concave
**

function.

**• On R, a continuous function is convex if and only if for all x, y ∈ R, we have
**

x+y f (x) + f (y)

f ≤ .

2 2

**• A twice continuously differentiable function on R is convex if and only if f 00 (x) ≥ 0 for
**

all x ∈ R.

• A local minimum of a convex function is a global minimum. See this page.

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Examples

**• ex ,e−x , and x2 are convex functions on R.
**

• A norm is a convex function.

• On R2 , the 1-norm and the ∞-norm (i.e., ||(x, y)||1 = |x| + |y| and ||(x, y)||∞ =

max{|x|, |y|}) are not strictly convex ([2], pp. 334-335).

REFERENCES

1. E. Kreyszig, Introductory Functional Analysis With Applications, John Wiley & Sons,

1978.

Version: 11 Owner: matte Author(s): matte, drini

382.2 extremal value of convex/concave functions

**Theorem. Let U be a convex set in a normed (real or complex) vector space. If f : U → R
**

is a convex function on U, then a local minimum of f is a global minimum.

**Proof. Suppose x is a local minimum for f , i.e., there is an open ball B ⊂ U with radius
**

and center x such that f (x) ≤ f (ξ) for all ξ ∈ B. Let us fix some y ∈ / B. Our aim is to

1

prove that f (x) ≤ f (y). We define λ = 2||x−y|| , where || · || is the norm on U.

Then

||λy + (1 − λ)x − x|| = ||λy − λx||

= |λ|||x − y||

= ,

2

so λy + (1 − λ)x ∈ B. If follows that f (x) ≤ f (λy + (1 − λ)x). Since f is convex, we then

get

f (x) ≤ f (λy + (1 − λ)x)

≤ λf (y) + (1 − λ)f (x),

and f (x) ≤ f (y) as claimed. 2

The analogous theorem for concave functions is as follows.

**Theorem. Let U be a convex set in a normed (real or complex) vector space. If f : U → R
**

is a concave function on U, then a local maximum of f is a global maximum.

1500

Proof. Consider the convex function −f . If x is a local maximum of f , then it is a local

minimum of −f . By the previous theorem, x is then a global minimum of −f . Hence x is a

global maximum of f . 2

Version: 1 Owner: matte Author(s): matte

1501

Chapter 383

**26B30 – Absolutely continuous
**

functions, functions of bounded

variation

383.1 absolutely continuous function

**Definition [1, 1] closed bounded interval of R. Then a function f : [a, b] → C is absolutely
**

continuous on [a, b], if for any ε > 0, there is a δ > 0 such that the following condition

holds:

**(∗) If (a1 , b1 ), . . . , (an , bn ) is a finite collection of disjoint open intervals in [a, b] such that
**

n

X

(bi − ai ) < δ,

i=1

then n

X

|f (bi ) − f (ai )| < ε.

i=1

Basic results for absolutely continuous functions are as follows.

Theorem

**1. A function f : [a, b] → C is absolutely continuous if and only if Re{f } and Im{f } are
**

absolutely continuous real functions.

**2. If f : [a, b] → C is a function, which is everywhere differentiable and f 0 is bounded,
**

then f is absolutely continuous [1].

1502

3. Any absolutely continuous function f : [a, b] → C is continuous on [a, b] and has a

bounded variation [1, 1].

**4. If f, g be absolutely continuous functions, then so are f g, f + g, |f |γ (if γ ≥ 1), and
**

f /g (if g is never zero) [1].

**5. If f, g are real valued absolutely continuous functions, then so are max{f, g} and
**

min{f, g}. If f (x) > 0 for all x and γ ∈ R, then f γ is absolutely continuous [1].

Property (2), which is readily proven using the mean value theorem, implies that any smooth function

with compact support on R is absolutely continuous. By property (3), any absolutely con-

tinuous function is a bounded variation. Hence, from properties of functions of bounded

variation, the following theorem follows:

**Theorem ([1], pp. 536) Let f : [a, b] → C be a absolutely continuous function. Then f is
**

differentiable almost everywhere, and |f 0| is Lebesgue integrable.

We have the following characterization of absolutely continuous functions

**Theorem [Fundamental theorem of calculus for the Lebesgue integral] ([1], pp. 550, [1])
**

Let f : [a, b] → C be a function. Then f is absolutely continuous if and only if there is a

function g ∈ L1 (a, b) (i.e. a g : (a, b) → C with int(a,b) |g| < ∞), such that

f (x) = f (a) + intxa g(t)dt

**for all x ∈ [a, b]. What is more, if f and g are as above, then f 0 = g almost everywhere.
**

(Above, both integrals are Lebesgue integrals.)

REFERENCES

1. G.B. Folland, Real Analysis: Modern Techniques and Their Applications, 2nd ed,

John Wiley & Sons, Inc., 1999.

2. W. Rudin, Real and complex analysis, 3rd ed., McGraw-Hill Inc., 1987.

3. F. Jones, Lebesgue Integration on Euclidean Spaces, Jones and Barlett Publishers, 1993.

4. C.D. Aliprantis, O. Burkinshaw, Principles of Real Analysis, 2nd ed., Academic Press,

1990.

Version: 5 Owner: matte Author(s): matte

383.2 total variation

**Let γ : [a, b] → X be a function mapping an interval [a, b] to a metric space (X, d). We
**

say that γ is of bounded variation if there is a constant M such that, for each partition

1503

P = {a = t0 < t1 < · · · < tn = b} of [a, b],

n

X

v(γ, P ) = d(γ(tk ), γ(tk−1 )) 6 M.

k=1

The total variation Vγ of γ is defined by

Vγ = sup{v(γ, P ) : P is a partition of [a, b]}.

**It can be shown that, if X is either R or C, every smooth (or piecewise smooth) function
**

γ : [a, b] → X is of bounded variation, and

Vγ = intba |γ 0 (t)|dt.

**Also, if γ is of bounded variation and f : [a, b] → X is continuous, then the Riemann-Stieltjes integral
**

intba f dγ is finite.

If γ is also continuous, it is said to be a rectifiable path, and V (γ) is the length of its trace.

**If X = R, it can be shown that γ is of bounded variation if and only if it is the difference of
**

two monotonic functions.

Version: 3 Owner: Koro Author(s): Koro

1504

Chapter 384

26B99 – Miscellaneous

384.1 derivation of zeroth weighted power mean

**Let x1 , x2 , . . . , xn be positive real numbers, and let w1 , w2 , . . . , wn be positive real numbers
**

such that w1 + w2 + · · · + wn = 1. For r 6= 0, the r-th weighted power mean of x1 , x2 , . . . , xn

is

Mwr (x1 , x2 , . . . , xn ) = (w1 xr1 + w2 xr2 + · · · + wn xrn )1/r .

Using the Taylor series expansion et = 1 + t + O(t2 ), where O(t2 ) is Landau notation for

terms of order t2 and higher, we can write xri as

xri = er log xi = 1 + r log xi + O(r 2 ).

**By substituting this into the definition of Mwr , we get
**

1/r

Mwr (x1 , x2 , . . . , xn ) = w1 (1 + r log x1 ) + · · · + wn (1 + r log xn ) + O(r 2 )

1/r

= 1 + r(w1 log x1 + · · · + wn log xn ) + O(r 2 )

2 1/r

= 1 + r log(xw 1 w2 wn

1 x2 · · · xn ) + O(r )

1 w1 w2 wn 2

= exp log 1 + r log(x1 x2 · · · xn ) + O(r ) .

r

Again using a Taylor series, this time log(1 + t) = t + O(t2 ), we get

r 1 w1 w2 wn 2

Mw (x1 , x2 , . . . , xn ) = exp r log(x1 x2 · · · xn ) + O(r )

r

= exp [log(xw 1 w2 wn

1 x2 · · · xn ) + O(r)] .

Taking the limit r → 0, we find

**Mw0 (x1 , x2 , . . . , xn ) = exp [log(xw 1 w2 wn
**

1 x2 · · · xn )]

= xw 1 w2 wn

1 x2 · · · xn .

1505

In particular, if we choose all the weights to be n1 ,

√

M 0 (x1 , x2 , . . . , xn ) = n

x1 x2 · · · xn ,

the geometric mean of x1 , x2 , . . . , xn .

Version: 3 Owner: pbruin Author(s): pbruin

384.2 weighted power mean

**If w1 , w2 , . . . , wn are positive real numbers such that w1 + w2 + · · · + wn = 1, we define the
**

r-th weighted power mean of the xi as:

Mwr (x1 , x2 , . . . , xn ) = (w1 xr1 + w2 xr2 + · · · + wn xrn )1/r .

**When all the wi = n1 we get the standard power mean. The weighted power mean is a
**

continuous function of r, and taking limit when r → 0 gives us

Mw0 = xw 1 w2 wn

1 x2 · · · wn .

**We can weighted use power means to generalize the power means inequality: If w is a set of
**

weights, and if r < s then

Mwr < Mws .

Version: 6 Owner: drini Author(s): drini

1506

Chapter 385

26C15 – Rational functions

385.1 rational function

A real function R(x) of a single variable x is called rational if it can be written as a quotient

P (x)

R(x) = ,

Q(x)

where P (x) and Q(x) are polynomials in x with real coefficients.

In general, a rational function R(x1 , . . . , xn ) has the form

P (x1 , . . . , xn )

R(x1 , . . . , xn ) = ,

Q(x1 , . . . , xn )

**where P (x1 , . . . , xn ) and Q(x1 , . . . , xn ) are polynomials in the variables (x1 , . . . , xn ) with
**

coefficients in some field or ring S.

**In this sense, R(x1 , . . . , xn ) can be regarded as an element of the fraction field S(x1 , . . . , xn )
**

of the polynomial ring S[x1 , . . . , xn ].

Version: 1 Owner: igor Author(s): igor

1507

Chapter 386

26C99 – Miscellaneous

386.1 Laguerre Polynomial

**A Laguerre Polynomial is a polynomial of the form:
**

ex dn −x n

Ln (x) = e x .

n! dxn

Associated to this is the Laguerre differential equation, the solutions of which are called

associated Laguerre Polynomials:

ex x−k dn −x n+k

Lkn (x) = e x .

n! dxn

Of course

L0n (x) = Ln (x).

The associated Laguere Polynomials are orthogonal over |0, ∞) with respect to the weighting

function xk e−x :

(n + k)!

int∞ x k k k

0 e x Ln (x)Lm (x)dx = δn m.

n!

Version: 2 Owner: mathwizard Author(s): mathwizard

1508

Chapter 387

**26D05 – Inequalities for trigonometric
**

functions and polynomials

387.1 Weierstrass product inequality

**For any finite family (ai )i∈I of real numbers in the interval [0, 1], we have
**

Y X

(1 − ai ) ≥ 1 − ai .

i i

Proof: Write Y X

f= (1 − ai ) + ai .

i i

**For any k ∈ I, and any fixed values of the ai for i 6= k, f is a polynomial of the first degree
**

in ak . Consequently f is minimal either at ak = 0 or ak = 1. That brings us down to two

cases: all the ai are zero, or at least one of them is 1. But in both cases it is clear that f ≥ 1,

QED.

Version: 2 Owner: Daume Author(s): Larry Hammick

387.2 proof of Jordan’s Inequality

**To prove that
**

2 π

x 6 sin(x) 6 x ∀ x ∈ [0, ] (387.2.1)

π 2

consider a unit circle (circle with radius = 1 unit). Take any point P on the circumference

of the circle.

1509

Drop the perpendicular from P to the horizontal line, M being the foot of the perpendicular

and Q the reflection of P at M. (refer to figure)

Let x = ∠P OM.

For x to be in [0, π2 ], the point P lies in the first quadrant, as shown.

**The length of line segment P M is sin(x). Construct a circle of radius MP , with M as the
**

center.

Length of line segment P Q is 2 sin(x).

Length of arc P AQ is 2x.

Length of arc P BQ is π sin(x).

**Since P Q 6 length of arc P AQ (equality holds when x = 0) we have 2 sin(x) 6 2x. This
**

implies

sin(x) 6 x

**Since length of arc P AQ is 6 length of arc P BQ (equality holds true when x = 0 or x = π2 ),
**

we have 2x 6 π sin(x). This implies

2

x 6 sin(x)

π

Thus we have

2 π

x 6 sin(x) 6 x ∀ x ∈ [0, ] (387.2.2)

π 2

Version: 12 Owner: giri Author(s): giri

1510

Chapter 388

**26D10 – Inequalities involving
**

derivatives and differential and

integral operators

388.1 Gronwall’s lemma

**If, for t0 6 t 6 t1 , φ(t) > 0 and ψ(t) > 0 are continuous functions such that the inequality
**

φ(t) 6 K + Linttt0 ψ(s)φ(s)ds

holds on t0 6 t 6 t1 , with K and L positive constants, then

φ(t) 6 K exp Linttt0 ψ(s)ds

on t0 6 t 6 t1 .

Version: 1 Owner: jarino Author(s): jarino

388.2 proof of Gronwall’s lemma

The inequality

φ(t) 6 K + Linttt0 ψ(s)φ(s)ds (388.2.1)

is equivalent to

φ(t)

61

K + Linttt0 ψ(s)φ(s)ds

Multiply by Lψ(t) and integrate, giving

Lψ(s)φ(s)ds

inttt0 6 Linttt0 ψ(s)ds

K + Lintst0 ψ(τ )φ(τ )dτ

1511

Thus

ln K + Linttt0 ψ(s)φ(s)ds − ln K 6 Linttt0 ψ(s)ds

and finally

K + Linttt0 ψ(s)φ(s)ds 6 K exp Linttt0 ψ(s)ds

Using (387.2.1) in the left hand side of this inequality gives the result.

Version: 2 Owner: jarino Author(s): jarino

1512

Chapter 389

**26D15 – Inequalities for sums, series
**

and integrals

389.1 Carleman’s inequality

**Theorem ([4], pp. 24) For positive real numbers {an }∞ n=1 , Carleman’s inequality states
**

that ∞ ∞

X 1/n X

a1 a2 · · · an ≤e an .

n=1 n=1

**Although the constant e (the natural log base) is optimal, it is possible to refine Carleman’s
**

inequality by decreasing the weight coefficients on the right hand side [2].

REFERENCES

1. L. Hörmander, The Analysis of Linear Partial Differential Operators I, (Distribution

theory and Fourier Analysis), 2nd ed, Springer-Verlag, 1990.

2. B.Q. Yuan, Refinements of Carleman’s inequality, Journal of Inequalities in Pure and

Applied Mathematics, Vol. 2, Issue 2, 2001, Article 21. online

Version: 2 Owner: matte Author(s): matte

389.2 Chebyshev’s inequality

**If x1 , x2 , . . . , xn and y1 , y2, . . . , yn are two sequences (at least one of them consisting of positive
**

numbers):

1513

• if x1 < x2 < · · · < xn and y1 < y2 < · · · < yn then

x1 + x2 + · · · + xn y1 + y2 + · · · + yn x1 y1 + x2 y2 + · · · + xn yn

≤ .

n n n

• if x1 < x2 < · · · < xn and y1 > y2 > · · · > yn then

x1 + x2 + · · · + xn y1 + y2 + · · · + yn x1 y1 + x2 y2 + · · · + xn yn

≥ .

n n n

Version: 1 Owner: drini Author(s): drini

389.3 MacLaurin’s Inequality

**Let a1 , a2 , . . . , an be positive real numbers , and define the sums Sk as follows :
**

X

ai1 ai2 · · · aik

16i1 <i2 <···<ik 6n

Sk =

n

k

Then the following chain of inequalities is true :

p p p

S1 > S2 > 3 S3 > · · · > n Sn

Note : Sk are called the averages of the elementary symmetric sums

This inequality is in fact important because it shows that the Arithmetic-Geometric Mean

inequality is nothing but a consequence of a chain of stronger inequalities

Version: 2 Owner: drini Author(s): drini, slash

389.4 Minkowski inequality

**If p > 1 and ak , bk are real numbers for k = 1, . . ., then
**

n

!1/p n

!1/p n

!1/p

X X X

|ak + bk |p ≤ |ak |p + |bk |p

k=1 k=1 k=1

**The Minkowski inequality is in fact valid for all Lp norms with p ≥ 1 on arbitrary measure spaces.
**

This covers the case of Rn listed here as well as spaces of sequences and spaces of functions,

and also complex Lp spaces.

Version: 8 Owner: drini Author(s): drini, saforres

1514

389.5 Muirhead’s theorem

**Let 0 6 s1 6 · · · 6 sn and 0 6 t1 6 . . . 6 tn be real numbers such that
**

n

X n

X k

X k

X

si = ti and si 6 ti (k = 1, . . . , n − 1)

i=1 i=1 i=1 i=1

**Then for any nonnegative numbers x1 , . . . , xn ,
**

X sσ(1) s X tσ(1) t

x1 . . . xnσ(n) 6 x1 . . . xnσ(n)

σ σ

where the sums run over all permutations σ of {1, 2, . . . , n}.

Version: 3 Owner: Koro Author(s): Koro

389.6 Schur’s inequality

**If a, b, and c are positive real numbers and k > 1 a fixed real constant, then the following
**

inequality holds:

ak (a − b)(b − c) + bk (b − c)(c − a) + ck (c − a)(c − b) > 0

Taking k = 1, we get the well-known

a3 + b3 + c3 + 3abc > ab(a + b) + ac(a + c) + bc(b + c)

**W e can assume without loss of generality that c 6 b 6 a via a permutation of the variables
**

(as both sides are symmetric in those variables). Then collecting terms, the lemma states

that

(a − b) ak (a − c) − bk (b − c) + ck (a − c)(b − c) > 0

which is clearly true as every term on the left is positive.

Version: 3 Owner: mathcam Author(s): mathcam, slash

389.7 Young’s inequality

**Let φ : R → R be a continuous , strictly increasing function such that φ(0) = 0 . Then the
**

following inequality holds:

ab 6 inta0 φ(x)dx + intb0 φ−1 (y)dy

1515

The inequality is trivial to prove by drawing the graph of φ(x) and by observing that the

sum of the two areas represented by the integrals above is greater than the area of a rectangle

of sides a and b .

Version: 2 Owner: slash Author(s): slash

389.8 arithmetic-geometric-harmonic means inequality

**Let x1 , x2 , . . . , xn be positive numbers. Then
**

x1 + x2 + · · · + xn

max{x1 , x2 , . . . , xn } ≥

√ n

≥ n x1 x2 · · · xn

n

≥ 1

x1

+ x2 + · · · + x1n

1

≥ min{x1 , x2 , . . . , xn }

There are several generalizations to this inequality using power means and weighted power means.

Version: 4 Owner: drini Author(s): drini

389.9 general means inequality

The power means inequality is a generalization of arithmetic-geometric means inequality.

**If 0 6= r ∈ R, the r-mean (or r-th power mean) of the nonnegative numbers a1 , . . . , an is
**

defined as !1/r

n

X

1

M r (a1 , a2 , . . . , an ) = ark

n

k=1

Given real numbers x, y such that xy 6= 0 and x < y, we have

Mx 6 My

and the equality holds if and only if a1 = ... = an .

**Additionally, if we define M 0 to be the geometric mean (a1 a2 ...an )1/n , we have that the
**

inequality above holds for arbitrary real numbers x < y.

**The mentioned inequality is a special case of this one, since M 1 is the arithmetic mean, M 0
**

is the geometric mean and M −1 is the harmonic mean.

1516

This inequality can be further generalized using weighted power means.

Version: 3 Owner: drini Author(s): drini

389.10 power mean

The r-th power mean of the numbers x1 , x2 , . . . , xn is defined as:

1/r

r xr1 + xr2 + · · · + xrn

M (x1 , x2 , . . . , xn ) = .

n

**The arithmetic mean is a special case when r = 1. The power mean is a continuous function
**

of r, and taking limit when r → 0 gives us the geometric mean:

√

M 0 (x1 , x2 , . . . , xn ) = n

x1 x2 · · · xn .

**Also, when r = −1 we get
**

n

M −1 (x1 , x2 , . . . , xn ) = 1 1 1

x1

+ x2

+···+ xn

the harmonic mean.

A generalization of power means are weighted power means.

Version: 8 Owner: drini Author(s): drini

389.11 proof of Chebyshev’s inequality

**Let x1 , x2 , . . . , xn and y1 , y2, . . . , yn be real numbers such that x1 ≤ x2 ≤ · · · ≤ xn . Write
**

the product (x1 + x2 + · · · + xn )(y1 + y2 + · · · + yn ) as

(x1 y1 + x2 y2 + · · · + xn yn )

+ (x1 y2 + x2 y3 + · · · + xn−1 yn + xn y1 )

+ (x1 y3 + x2 y4 + · · · + xn−2 yn + xn−1 y1 + xn y2 )

+ ···

+ (x1 yn + x2 y1 + x3 y2 + · · · + xn yn−1 ). (389.11.1)

1517

• If y1 ≤ y2 ≤ · · · ≤ yn , each of the n terms in parentheses is less than or equal to

x1 y1 + x2 y2 + · · · + xn yn , according to the rearrangement inequality. From this, it

follows that

(x1 + x2 + · · · + xn )(y1 + y2 + · · · + yn ) ≤ n(x1 y1 + x2 y2 + · · · + xn yn )

or (dividing by n2 )

x1 + x2 + · · · + xn y1 + y2 + · · · + yn x1 y1 + x2 y2 + · · · + xn yn

≤ .

n n n

**• If y1 ≥ y2 ≥ · · · ≥ yn , the same reasoning gives
**

x1 + x2 + · · · + xn y1 + y2 + · · · + yn x1 y1 + x2 y2 + · · · + xn yn

≥ .

n n n

**It is clear that equality holds if x1 = x2 = · · · = xn or y1 = y2 = · · · = yn . To see that
**

this condition is also necessary, suppose that not all yi ’s are equal, so that y1 6= yn . Then

the second term in parentheses of (388.11.1) can only be equal to x1 y1 + x2 y2 + · · · + xn yn

if xn−1 = xn , the third term only if xn−2 = xn−1 , and so on, until the last term which can

only be equal to x1 y1 + x2 y2 + · · · + xn yn if x1 = x2 . This implies that x1 = x2 = · · · = xn .

Therefore, Chebyshev’s inequality is an equality if and only if x1 = x2 = · · · = xn or

y1 = y2 = · · · = yn .

Version: 1 Owner: pbruin Author(s): pbruin

389.12 proof of Minkowski inequality

**For p = 1 the result follows immediately from the triangle inequality, so we may assume
**

p > 1.

We have

|ak + bk |p = |ak + bk ||ak + bk |p−1 6 (|ak | + |bk |)|ak + bk |p−1

by the triangle inequality. Therefore we have

|ak + bk |p 6 |ak ||ak + bk |p−1 + |bk ||ak + bk |p−1

p 1 1

Set q = p−1

. Then p

+ q

= 1, so by the Hölder inequality we have

n n

! 1p n

! 1q

X X X

|ak ||ak + bk |p−1 6 |ak |p |ak + bk |(p−1)q

k=0 k=0 k=0

n n

! p1 n

! 1q

X X X

|bk ||ak + bk |p−1 6 |bk |p |ak + bk |(p−1)q

k=0 k=0 k=0

1518

Adding these two inequalities, dividing by the factor common to the right sides of both, and

observing that (p − 1)q = p by definition, we have

n

!1− 1q n n

! p1 n

! p1

X X X X

|ak + bk |p 6 (|ak | + |bk |)|ak + bk |p−1 6 |ak |p + |bk |p

i=0 k=0 k=0 k=0

**Finally, observe that 1 − 1q = 1p , and the result follows as required. The proof for the integral
**

version is analogous.

Version: 4 Owner: saforres Author(s): saforres

**389.13 proof of arithmetic-geometric-harmonic means
**

inequality

Let M be max{x1 , x2 , x3 , . . . , xn } and let m be min{x1 , x2 , x3 , . . . , xn }.

Then

M +M +M +···+M x1 + x2 + x3 + · · · + xn

M= >

n n

n n n

m= n = 1 1 1 1 6 1

m m

+ m + m +···+ m x1

+ x2 + x3 + · · · + x1n

1 1

**where all the summations have n terms. So we have proved in this way the two inequalities
**

at the extremes.

**Now we shall prove the inequality between arithmetic mean and geometric mean. We do
**

first the case n = 2.

√ √

( x1 − x2 )2 > 0

√

x1 − 2 x1 x2 + x2 > 0

√

x1 + x2 > 2 x1 x2

x1 + x2 √

> x1 x2

2

**Now we prove the inequality for any power of 2 (that is, n = 2k for some integer k) by using
**

mathematical induction.

**x1 + x2 + · · · + x2k + x2k +1 + · · · + x2k+1
**

2k+1

x +x +···+x2k+1

x1 +x2 +···+x2k 2k +1 2k +2

2k + 2k

=

2

1519

and using the case n = 2 on the last expression we can state the following inequality

x1 + x2 + · · · + x2k + x2k +1 + · · · + x2k+1

s 2k+1

x1 + x2 + · · · + x2k x2k +1 + x2k +2 + · · · + x2k+1

≥

2k 2k

q

√

2k

√

≥ x1 x2 · · · x2k 2k x2k +1 x2k +2 · · · x2k+1

**where the last inequality was obtained by applying the induction hypothesis with n = 2k .
**

√

Finally, we see that the last expression is equal to 2k+1 x1 x2 x3 · · · x2k+1 and so we have proved

the truth of the inequality when the number of terms is a power of two.

**Finally, we prove that if the inequality holds for any n, it must also hold for n − 1, and
**

this proposition, combined with the preceding proof for powers of 2, is enough to prove the

inequality for any positive integer.

Suppose that

x1 + x2 + · · · + xn √

> n x1 x2 · · · xn

n

is known for a given value of n (we just proved that it is true for powers of two, as example).

Then we can replace xn with the average of the first n − 1 numbers. So

+···+xn−1

x1 + x2 + · · · + xn−1 + x1 +x2n−1

n

(n − 1)x1 + (n − 1)x2 + · · · + (n − 1)xn−1 + x1 + x2 + · · · + xn

=

n(n − 1)

nx1 + nx2 + · · · + nxn−1

=

n(n − 1)

x1 + x2 + · · · + xn−1

=

(n − 1)

**On the other hand
**

s

n x1 + x2 + · · · + xn−1

x1 x2 · · · xn−1

n−1

r

√ x1 + x2 + · · · + xn−1

= n x1 x2 · · · xn−1 n

n−1

which, by the inequality stated for n and the observations made above, leads to:

n

x1 + x2 + · · · + xn−1 x1 + x2 + · · · + xn−1

≥ (x1 x2 · · · xn )

n−1 n−1

and so n−1

x1 + x2 + · · · + xn−1

≥ x1 x2 · · · xn

n−1

1520

from where we get that

x1 + x2 + · · · + xn−1 √

≥ n−1

x1 x2 · · · xn .

n−1

So far we have proved the inequality between the arithmetic mean and the geometric mean.

The geometric-harmonic inequality is easier. Let ti be 1/xi .

From

t1 + t2 + · · · + tn √

> n t1 t2 t3 · · · tn

n

we obtain r

1

x1

+ x12 + x13 + · · · + x1n 1 1 1 1

> n ···

n x1 x2 x3 xn

and therefore

√ n

n

x1 x2 x3 · · · xn > 1 1 1 1

x1

+ x2

+ x3

+···+ xn

and so, our proof is completed.

Version: 2 Owner: drini Author(s): drini

389.14 proof of general means inequality

**Let r < s be real numbers, and let w1 , w2 , . . . , wn be positive real numbers such that w1 +
**

w2 + · · · + wn = 1. We will prove the weighted power means inequality, which states that

for positive real numbers x1 , x2 , . . . , xn ,

Mwr (x1 , x2 , . . . , xn ) ≤ Mws (x1 , x2 , . . . , xn ).

**First, suppose that r and s are nonzero. Then the r-th weighted power mean of x1 , x2 , . . . , xn
**

is

Mwr (x1 , x2 , . . . , xn ) = (w1 x1 + w2 x2 + · · · + wn xn )1/r

and Mws is defined similarly.

**Let t = rs , and let yi = xri for 1 ≤ i ≤ n; this implies yit = xsi . Define the function f on
**

1

(0, ∞) by f (x) = xt . The second derivative of f is f 00 (x) = t(t−1) xt−2 . There are three cases

for the signs of r and s: r < s < 0, r < 0 < s, and 0 < r < s. We will prove the inequality

for the case 0 < r < s; the other cases are almost identical.

1

In the case that r and s are both positive, t > 1. Since f 00 (x) = t(t−1) xt−2 > 0 for all x > 0,

f is a strictly convex function. Therefore, according to Jensen’s inequality,

(w1 y1 + w2 y2 + · · · + wn yn )t = f (w1y1 + w2 y2 + · · · + wn yn )

≤ w1 f (y1) + w2 f (y2) + · · · + wn f (yn )

= w1 y1t + w2 y2t + · · · + wn ynt .

1521

s

with equality if and only if y1 = y2 = · · · = yn . By substituting t = r

and yi = xri back into

this inequality, we get

(w1 xr1 + w2 xr2 + · · · + wn xrn )s/r ≤ w1 xs1 + w2 xs2 + · · · + wn xsn

**with equality if and only if x1 = x2 = · · · = xn . Since s is positive, the function x 7→ x1/s is
**

strictly increasing, so raising both sides to the power 1/s preserves the inequality:

(w1 xr1 + w2 xr2 + · · · + wn xrn )1/r ≤ (w1 xs1 + w2 xs2 + · · · + wn xsn )1/s ,

which is the inequality we had to prove. Equality holds if and only if all the xi are equal.

**If r = 0, the inequality is still correct: Mw0 is defined as limr→0 Mwr , and since Mwr ≤ Mws
**

for all r < s with r 6= 0, the same holds for the limit r → 0. We can show by an identical

argument that Mwr ≤ Mw0 for all r < 0. Therefore, for all real numbers r and s such that

r < s,

Mwr (x1 , x2 , . . . , xn ) ≤ Mws (x1 , x2 , . . . , xn ).

Version: 1 Owner: pbruin Author(s): pbruin

389.15 proof of rearrangement inequality

**We first prove the rearrangement inequality for the case n = 2. Let x1 , x2 , y1 , y2 be real numbers
**

such that x1 ≤ x2 and y1 ≤ y2 . Then

(x2 − x1 )(y2 − y1 ) ≥ 0,

and therefore

x1 y1 + x2 y2 ≥ x1 y2 + x2 y1 .

Equality holds iff x1 = x2 or y1 = y2 .

**For the general case, let x1 , x2 , . . . , xn and y1 , y2, . . . , yn be real numbers such that x1 ≤ x2 ≤
**

· · · ≤ xn . Suppose that (z1 , z2 , . . . , zn ) is a permutation (rearrangement) of {y1, y2 , . . . , yn }

such that the sum

x1 z1 + x2 z2 + · · · + xn zn

is maximized. If there exists a pair i < j with zi > zj , then xi zj + xj zi ≥ xi zi + xj zj (the

n = 2 case); equality holds iff xi = xj . Therefore, x1 z1 + x2 z2 + · · · + xn zn is not maximal

unless z1 ≤ z2 ≤ · · · ≤ zn or xi = xj for all pairs i < j such that zi > zj . In the latter case, we

can consecutively interchange these pairs until z1 ≤ z2 ≤ · · · ≤ zn (this is possible because

the number of pairs i < j with zi > zj decreases with each step). So x1 z1 + x2 z2 + · · · + xn zn

is maximized if

z1 ≤ z2 ≤ · · · ≤ zn .

To show that x1 z1 + x2 z2 + · · · + xn zn is minimal for a permutation (z1 , z2 , . . . , zn ) of

{y1 , y2 , . . . , yn } if z1 ≥ z2 ≥ · · · ≥ zn , observe that −(x1 z1 + x2 z2 + · · · + xn zn ) = x1 (−z1 ) +

1522

x2 (−z2 ) + · · · + xn (−zn ) is maximized if −z1 ≤ −z2 ≤ · · · ≤ −zn . This implies that

x1 z1 + x2 z2 + · · · + xn zn is minimized if

z1 ≥ z2 ≥ · · · ≥ zn .

Version: 1 Owner: pbruin Author(s): pbruin

389.16 rearrangement inequality

Let x1 , x2 , . . . , xn and y1 , y2 , . . . , yn two sequences of positive real numbers. Then the sum

x1 y1 + x2 y2 + · · · + xn yn

**is maximized when the two sequences are ordered in the same way (i.e. x1 ≤ x2 ≤ · · · ≤ xn
**

and y1 ≤ y2 ≤ · · · ≤ yn ) and is minimized when the two sequences are ordered in the opposite

way (i.e. x1 ≤ x2 ≤ · · · ≤ xn and y1 ≥ y2 ≥ · · · ≥ yn ).

**This can be seen intuitively as: If x1 , x2 , . . . , xn are the prices of n kinds of items, and
**

y1 , y2 , . . . , yn the number of units sold of each, then the highest profit is when you sell more

items with high prices and fewer items with low prices (same ordering), and the lowest profit

happens when you sell more items with lower prices and less items with high prices (opposite

orders).

Version: 4 Owner: drini Author(s): drini

1523

Chapter 390

26D99 – Miscellaneous

390.1 Bernoulli’s inequality

Let x and r be real numbers. If r > 1 and x > −1 then

(1 + x)r ≥ 1 + xr.

The inequality also holds when r is an even integer.

Version: 3 Owner: drini Author(s): drini

390.2 proof of Bernoulli’s inequality

Let I be the interval (−1, ∞) and f : I → R the function defined as:

f (x) = (1 + x)α − 1 − αx

with α ∈ R \ {0, 1} fixed. Then f is differentiable and its derivative is

f 0 (x) = α(1 + x)α−1 − α, for all x ∈ I,

from which it follows that f 0 (x) = 0 ⇔ x = 0.

**1. If 0 < α < 1 then f 0 (x) < 0 for all x ∈ (0, ∞) and f 0 (x) > 0 for all x ∈ (−1, 0)
**

which means that 0 is a global maximum point for f . Therefore f (x) < f (0) for all

x ∈ I \ {0} which means that (1 + x)α < 1 + αx for all x ∈ (−1, 0).

1524

/ [0, 1] then f 0 (x) > 0 for all x ∈ (0, ∞) and f 0 (x) < 0 for all x ∈ (−1, 0) meaning

2. If α ∈

that 0 is a global minimum point for f . This implies that f (x) > f (0) for all x ∈ I \{0}

which means that (1 + x)α > 1 + αx for all x ∈ (−1, 0).

Checking that the equality is satisfied for x = 0 or for α ∈ {0, 1} ends the proof.

Version: 3 Owner: danielm Author(s): danielm

1525

Chapter 391

26E35 – Nonstandard analysis

391.1 hyperreal

An ultrafilter F on a set I is called nonprincipal if no finite subsets of I are in F.

**Fix once and for all a nonprincipal ultrafilter F on the set N of natural numbers. Let ∼ be
**

the equivalence relation on the set RN of sequences of real numbers given by

{an } ∼ {bn } ⇔ {n ∈ N | an = bn } ∈ F

**Let ∗ R be the set of equivalence classes of RN under the equivalence relation ∼. The set ∗ R
**

is called the set of hyperreals. It is a field under coordinatewise addition and multiplication:

**{an } + {bn } = {an + bn }
**

{an } · {bn } = {an · bn }

The field ∗ R is an ordered field under the ordering relation

{an } 6 {bn } ⇔ {n ∈ N | an 6 bn } ∈ F

**The real numbers embed into ∗ R by the map sending the real number x ∈ R to the equivalence
**

class of the constant sequence given by xn := x for all n. In what follows, we adopt the

convention of treating R as a subset of ∗ R under this embedding.

A hyperreal x ∈ ∗ R is:

• limited if a < x < b for some real numbers a, b ∈ R

• positive unlimited if x > a for all real numbers a ∈ R

• negative unlimited if x < a for all real numbers a ∈ R

1526

• unlimited if it is either positive unlimited or negative unlimited

• positive infinitesimal if 0 < x < a for all positive real numbers a ∈ R+

• negative infinitesimal if a < x < 0 for all negative real numbers a ∈ R−

• infinitesimal if it is either positive infinitesimal or negative infinitesimal

**For any subset A of R, the set ∗ A is defined to be the subset of ∗ R consisting of equivalence
**

classes of sequences {an } such that

{n ∈ N | an ∈ A} ∈ F.

**The sets ∗ N, ∗ Z, and ∗ Q are called hypernaturals, hyperintegers, and hyperrationals, respec-
**

tively. An element of ∗ N is also sometimes called hyperfinite.

Version: 1 Owner: djao Author(s): djao

391.2 e is not a quadratic irrational

**Looking at the Taylor series for ex , we see that
**

∞

X

x xk

e = .

k=0

k!

P∞ 1 −1

P∞ k 1

This converges for every x ∈ R, so e = k=0 k! and e = k=0 (−1) k! . Arguing by

contradiction, assume ae2 +be+c = 0 for integers a, b and c. That is the same as ae+b+ce−1 =

0.

Fix n > |a| + |c|, then a, c | n! and ∀k ≤ n, k! | n! . Consider

X∞ X∞

−1 1 1

0 = n!(ae + b + ce ) = an! + b + cn! (−1)k

k=0

k! k=0

k!

n

X ∞

X

n! n!

= b+ (a + c(−1)k ) + (a + c(−1)k )

k! k!

k=0 k=n+1

Since k! | n! for k ≤ n, the first two terms are integers. So the third term should be an

1527

integer. However,

X∞

n! X∞

n!

(a + c(−1)k ) ≤ (|a| + |c|)

k! k!

k=n+1 k=n+1

∞

X 1

= (|a| + |c|)

k=n+1

(n + 1)(n + 2) · · · k

∞

X

≤ (|a| + |c|) (n + 1)n−k

k=n+1

∞

X

= (|a| + |c|) (n + 1)−t

t=1

1

= (|a| + |c|)

n

is less than 1 by our assumption that n > |a|P+ |c|. Since there is only one integer which is

less than 1 in absolute value, this means that ∞ k 1

k=n+1 (a+ c(−1) ) k! = 0 for every sufficiently

large n which is not the case because

∞

X X∞

1

k 1 1

(a + c(−1) ) − (a + c(−1)k ) = (a + c(−1)k )

k=n+1

k! k=n+2 k! (n + 1)!

is not identically zero. The contradiction completes the proof.

Version: 6 Owner: thedagit Author(s): bbukh, thedagit

391.3 zero of a function

**Definition Suppose X is a set, and suppose f is a complex-valued function f : X → C.
**

Then a zero of f is an element x ∈ X such that f (x) = 0. The zero set of f is the set

Z(f ) = {x ∈ X | f (x) = 0}.

Remark

**When X is a “simple” space, such as R or C a zero is also called a root. However, in pure
**

mathematics and especially if Z(f ) is infinite, it seems to be customary to talk of zeroes and

the zero set instead of roots.

Examples

**• Suppose p is a polynomial p : C → C of degree n ≥ 1. Then p has at most n zeroes.
**

That is, |Z(p)| ≤ n.

1528

• If f and g are functions f : X → C and g : X → C, then

[

Z(f g) = Z(f ) Z(g),

Z(f g) ⊃ Z(f ),

where f g is the function x 7→ f (x)g(x).

• If X is a topological space and f : X → C is a function, then

supp f = Z(f ){.

**Further, if f is continuous, then Z(f ) is a closed in X (assuming that C is given the
**

usual topology of the complex plane where {0} is a closed set).

Version: 21 Owner: mathcam Author(s): matte, yark, say 10, apmxi

1529

Chapter 392

**28-00 – General reference works
**

(handbooks, dictionaries,

bibliographies, etc.)

392.1 extended real numbers

**The extended real numbers are the real numbers together with +∞ (or simply ∞) and
**

−∞. This set is usually denoted by R or [−∞, ∞] [3], and the elements +∞ and −∞ are

called plus infinity respectively minus infinity. Following [3], let us next extend the order

operation <, the addition and multiplication operations, and the absolute value from R to

R. In other words, let us define how these operations should behave when some of their

arguments are ∞ or −∞.

Order on R

The order relation on R extends to R by defining that for any x ∈ R, we have

−∞ < x,

x < ∞,

and that −∞ < ∞.

Addition

For any real number x, we define

x + (±∞) = (±∞) + x = ±∞,

1530

and for +∞ and −∞, we define

(±∞) + (±∞) = ±∞.

It should be pointed out that sums like (+∞) + (−∞) are left undefined.

Multiplication

If x is a positive real number, then

x · (±∞) = (±∞) · x = ±∞.

Similarly, if x is a negative real number, then

x · (±∞) = (±∞) · x = ∓∞.

Furthermore, for ∞ and −∞, we define

(+∞) · (+∞) = (−∞) · (−∞) = +∞,

(+∞) · (−∞) = (−∞) · (+∞) = −∞.

**In many areas of mathematics, products like 0 · ∞ are left undefined. However, a special
**

case is measure theory, where it is convenient to define [3]

0 · (±∞) = (±∞) · 0 = 0.

Absolute value

For ∞ and −∞, the absolute value is defined as

| ± ∞| = +∞.

Examples

1. By taking x = −1 in the product rule, we obtain the relations

(−1) · (±∞) = ∓∞.

REFERENCES

1. D.L. Cohn, Measure Theory, Birkhäuser, 1980.

Version: 1 Owner: matte Author(s): matte

1531

Chapter 393

28-XX – Measure and integration

393.1 Riemann integral

**Suppose there is a function f : D → R where D, R ⊆ R and that there is a closed interval I =
**

[a, b] such that I ⊆ D. For any finite set of points {x0 , x1 , x2 , . . . xn } such that a = x0 < x1 <

x2 · · · < xn = b, there is a corresponding partition P = {[x0 , x1 ), [x1 , x2 ), . . . , [xn − 1, xn ]} of

I.

**Let C() be the set of all partitions of I with max(xi+1 − xi ) < . Then let S ∗ () be the
**

infimum of the set of upper Riemann sums with each partition in C(), and let S∗ () be the

supremum of the set of lower Riemann sums with each partition in C(). If 1 < 2 , then

C(1 ) ⊂ C(2 ), so S ∗ = lim→0 S ∗ () and S∗ = lim→0 S∗ () exist. If S ∗ = S∗ , then f is

Riemann-integrable over I, and the Riemann integral of f over I is defined by

intba f (x)dx = S ∗ = S∗ .

Version: 4 Owner: bbukh Author(s): bbukh, vampyr

393.2 martingale

Let ν be a probability measure on Cantor space C, and let s ∈ [0, ∞).

**1. A ν-s-supergale is a function d : {0, 1}∗ → [0, ∞) that satisfies the condition
**

d(w)ν(w)s > d(w0)ν(w0)s + d(w1)ν(w1)s (393.2.1)

for all w ∈ {0, 1}∗.

2. A ν-s-gale is a ν-s-supergale that satisfies the condition with equality for all w ∈ {0, 1}∗.

1532

3. A ν-supermartingale is a ν-1-supergale.

4. A ν-martingale is a ν-1-gale.

5. An s-supergale is a µ-s-supergale, where µ is the uniform probability measure.

6. An s-gale is a µ-s-gale.

7. A supermartingale is a 1-supergale.

8. A martingale is a 1-gale.

**Put in another way, a martingale is a function d : {0, 1}∗ → [0, ∞) such that, for all
**

w ∈ {0, 1}∗, d(w) = (d(w0) + d(w1))/2.

**Let d be a ν-s-supergale, where ν is a probability measure on C and s ∈ [0, ∞). We say that
**

d succeeds on a sequence S ∈ C if

lim sup d(S[0..n − 1]) = ∞.

n→∞

The success set of d is S ∞ [d] = {S ∈ Cd succeeds on S}. d succeeds on a language A ⊆

{0, 1}∗ if d succeeds on the characteristic sequence χA of A. We say that d succeeds strongly

on a sequence S ∈ C if

lim inf d(S[0..n − 1]) = ∞.

n→∞

∞

The strong success set of d is Sstr [d] = {S ∈ Cd succeeds strongly on S}.

**Intuitively, a supergale d is a betting strategy that bets on the next bit of a sequence when
**

the previous bits are known. s is the parameter that tunes the fairness of the betting. The

smaller s is, the less fair the betting is. If d succeeds on a sequence, then the bonus we can

get from applying d as the betting strategy on the sequence is unbounded. If d succeeds

strongly on a sequence, then the bonus goes to infinity.

Version: 10 Owner: xiaoyanggu Author(s): xiaoyanggu

1533

Chapter 394

**28A05 – Classes of sets (Borel fields,
**

σ-rings, etc.), measurable sets, Suslin

sets, analytic sets

394.1 Borel σ-algebra

**For any topological space X, the Borel sigma algebra of X is the σ–algebra B generated by
**

the open sets of X. An element of B is called a Borel subset of X, or a Borel set.

Version: 5 Owner: djao Author(s): djao, rmilson

1534

Chapter 395

**28A10 – Real- or complex-valued set
**

functions

395.1 σ-finite

**A measure space (Ω, B, µ) is σ-finite if the total space is the union of a finite or countable
**

family of sets of finite measure; i.e. if there

S exists a finite or countable set F ⊂ B such that

µ(A) < ∞ for each A ∈ F, and Ω = A∈F A. In this case we also say that µ is a σ-finite

measure. If µ is not σ-finite, we say that it is σ-infinite.

**Examples. Any finite measure space is σ-finite. A more interesting example is the Lebesgue measure
**

µ in Rn : it is σ-finite but not finite. In fact

[

R= [−k, k]n

k∈N

**([−k, k]n is a cube with center at 0 and side length 2k, and its measure is (2k)n ), but
**

µ(Rn ) = ∞.

Version: 6 Owner: Koro Author(s): Koro, drummond

395.2 Argand diagram

**An argand diagram is the graphical representation of complex numbers written in polar
**

coordinates.

**Argand is the name of Jean-Robert Argand, the frenchman who is is credited with the
**

geometric interpretation of the complex numbers [Biography]

Version: 3 Owner: drini Author(s): drini

1535

395.3 Hahn-Kolmogorov theorem

S

Let A0 be an algebra of subsets of a set X. If a finitely additive measure µ0 : A → R {∞}

satisfies ∞ ∞

[ X

µ0 ( An ) = µ0 (An )

n=1 n=1

S

for any disjoint family {An : n ∈ N} of elements of A0 such that ∞ n=0 An ∈ A0 , then µ0

extends uniquely to a measureS defined on the σ-algebra A generated by A0 ; i.e. there exists

a unique measure µ : A → R {∞} such that its restriction to A0 coincides with µ0

Version: 3 Owner: Koro Author(s): Koro

395.4 measure

**S (E, B(E)) be a measurable space. A measure on (E, B(E)) is a function µ : B(E) −→
**

Let

R {∞} with values in the extended real numbers such that:

**1. µ(A) > 0 for A ∈ B(E), with equality if A = ∅
**

S P∞

2. µ( ∞i=0 Ai ) = i=0 µ(Ai ) for any sequence of disjoint sets Ai ∈ B(E).

I

**The second property is called countable additivity. A finitely additive measure µ has the
**

same definition except that B(E) is only required to be an algebra and the second property

above is only required to hold for finite unions. Note the slight abuse of terminology: a

finitely additive measure is not necessarily a measure.

**The triple (E, B, µ) is called a measure space. If µ(E) = 1, then it is called a probability
**

space, and the measure µ is called a probability measure.

Lebesgue measure on Rn is one important example of a measure.

Version: 8 Owner: djao Author(s): djao

395.5 outer measure

**Definition [1, 2, 1] Let X be a set, and let P(X) be the power set of X. An outer measure
**

on X is a function µ∗ : P(X) → [0, ∞] satisfying the properties

1. µ∗ (∅) = 0.

1536

2. If A ⊂ B are subsets in X, then µ∗ (A) ≤ µ∗ (B).

**3. If {Ai } is a countable collection of subsets of X, then
**

[ X

µ∗ ( Ai ) ≤ µ∗ (Ai ).

i i

**Here, we can make two remarks. First, from (1) and (2), it follows that µ∗ is a positive
**

function on P(X). Second, property (3) also holds for any finite collection of subsets since

we can always append an infinite sequence of empty sets to such a collection.

Examples

**• [1, 2] On a set X, let us define µ∗ : P(X) → [0, ∞] as
**

∗ 1 when E 6= ∅,

µ (E) =

0 when E = ∅.

Then µ∗ is an outer measure.

**• [1] On a uncountable set X, let us define µ∗ : P(X) → [0, ∞] as
**

∗ 1 when E is uncountable,

µ (E) =

0 when E is countable.

Then µ∗ is an outer measure.

**Theorem [1, 2, 1] Let X be a set, and let F be a collection of subsets of X such that ∅ ∈ F
**

and X ∈ F. Further, let ρ : F → [0, ∞] be a mapping such that ρ(∅) = 0. If A ⊂ X, let

∞

X

∗

µ (A) = inf ρ(Fi ),

i=1

S∞

where the infimum is taken over all collections {Fi }∞

i=1 ⊂ F such that A ⊂ i=1 Fi . Then

∗

µ : P(X) → [0, ∞] is an outer measure.

REFERENCES

1. A. Mukherjea, K. Pothoven, Real and Functional analysis, Plenum press, 1978.

2. A. Friedman, Foundations of Modern Analysis, Dover publications, 1982.

3. G.B. Folland, Real Analysis: Modern Techniques and Their Applications, 2nd ed,

John Wiley & Sons, Inc., 1999.

Version: 1 Owner: mathcam Author(s): matte

1537

395.6 properties for measure

**Theorem [1, 1, 3, 2] Let (E, B, µ) be a measure space, i.e., let E be a set, let B be a
**

σ-algebra of sets in E, and let µ be a measure on B. Then the following properties hold:

**1. Monotonicity: If A, B ∈ B, and A ⊂ B, then µ(A) ≤ µ(B).
**

2. If A, B in B, A ⊂ B, and µ(A) < ∞, then

µ(B \ A) = µ(B) − µ(A).

**3. For any A, B in B, we have
**

[ \

µ(A B) + µ(A B) = µ(A) + µ(B).

4. subadditivity: If {Ai }∞

i=1 is a collection of sets from B, then

∞

[ ∞

X

µ Ai ≤ µ(Ai ).

i=1 i=1

**5. Continuity from below: If {Ai }∞
**

i=1 is a collection of sets from B such that Ai ⊂ Ai+1

for all i, then

[∞

µ Ai = lim µ(Ai ).

i→∞

i=1

**6. Continuity from above: If {Ai }∞ i=1 is a collection of sets from B such that µ(A1 ) <
**

∞, and Ai ⊃ Ai+1 for all i, then

∞

\

µ Ai = lim µ(Ai ).

i→∞

i=1

**Remarks In (2), the assumption µ(A) < ∞ assures that the right hand side is always
**

well defined, i.e., not of the form ∞ − ∞. Without the assumption we can prove T that

µ(B) = µ(A) + µ(B \ A) (see below). In (3), it is tempting to move the term µ(A B) to

the other side for aesthetic reasons. However, this is only possible if the term is finite.

S

Proof. For (1), suppose A ⊂ B. We can then write B as the disjoint union B = A (B \ A),

whence

[

µ(B) = µ(A (B \ A)) = µ(A) + µ(B \ A).

**Since µ(B \ A) ≥ 0, the claim follows. Property (2) follows from the above equation; since
**

S < ∞, S

µ(A) we can subtract this quantity from both sides. For property (3), we can write

A B = A (B \ A), whence

[

µ(A B) = µ(A) + µ(B \ A)

≤ µ(A) + µ(B).

1538

S

If µ(A B) is infinite, the last inequality must be equality, and either of µ(A) or µ(B) T must

be infinite.

S Together with (1), we obtain that if any of the quantities µ(A), µ(B), µ(A B)

or µ(A B) is infinite, then all quantities are infinite, whence the claim clearly holds. S We

canS therefore without loss of generality assume that all quantities are finite. From A B=

B (A \ B), we have [

µ(A B) = µ(B) + µ(A \ B)

and thus [

2µ(A B) = µ(A) + µ(B) + µ(A \ B) + µ(B \ A).

For the last two terms we have

[

µ(A \ B) + µ(B \ A) = µ((A \ B) (B \ A))

[ \

= µ((A B) \ (A B))

[ \

= µ(A B) − µ(A B),

where, in the second equality we have used properties for the symmetric set difference, and

the last equality follows from property (2). This completes the proof of property (3). For

property (4), let us define the sequence {Di }∞

i=1 as

i−1

[

D1 = A1 , Di = Ai \ Ak .

k=1

T S S∞

Now Di Dj = ∅ for i < j, so {Di } is a sequence of disjoint sets. Since ∞i=1 Di = i=1 Ai ,

and since Di ⊂ Ai , we have

[∞ ∞

[

µ( Ai ) = µ( Di )

i=1 i=1

∞

X

= µ(Di )

i=1

X∞

≤ µ(Ai ),

i=1

and property (4) follows.

TODO: proofs for (5)-(6).

REFERENCES

1. G.B. Folland, Real Analysis: Modern Techniques and Their Applications, 2nd ed,

John Wiley & Sons, Inc., 1999.

2. A. Mukherjea, K. Pothoven, Real and Functional analysis, Plenum press, 1978.

3. D.L. Cohn, Measure Theory, Birkhäuser, 1980.

4. A. Friedman, Foundations of Modern Analysis, Dover publications, 1982.

Version: 2 Owner: matte Author(s): matte

1539

Chapter 396

**28A12 – Contents, measures, outer
**

measures, capacities

396.1 Hahn decomposition theorem

**Let µ be a signed measure in the measurable space (Ω, S). There are two measurable sets A
**

and B such that:

S T

1. A B = Ω and A B = ∅;

2. µ(E) > 0 for each E ∈ S such that E ⊂ A;

3. µ(E) 6 0 for each E ∈ S such that E ⊂ B.

**The pair (A, B) is called a Hahn decomposition for µ. This decomposition is not unique,
**

but any other such decomposition (A0 , B 0 ) satisfies µ(A0 M A) = µ(B M B 0 ) = 0 (where M

denotes the symmetric difference), so the two decompositions differ in a set of measure 0.

Version: 6 Owner: Koro Author(s): Koro

396.2 Jordan decomposition

**Let (Ω, S, µ) be a signed measure space, and let (A, B) be a Hahn decomposition for µ. We
**

define µ+ and µ− by

\ \

µ+ (E) = µ(A E) and µ− (E) = −µ(B E).

This definition is easily shown to be independent of the chosen Hahn decomposition.

1540

It is clear that µ+ is a positive measure, and it is called the positive variation of µ. On

the other hand, µ− is a positive finite measure, called the negative variation of µ. The

measure |µ| = µ+ + µ− is called the total variation of µ.

**Notice that µ = µ+ − µ− . This decomposition of µ into its positive and negative parts is
**

called the Jordan decomposition of µ.

Version: 6 Owner: Koro Author(s): Koro

396.3 Lebesgue decomposition theorem

**Let µ and ν be two σ-finite signed measures in the measurable space (Ω, S). There exist two
**

σ-finite signed measures ν0 and ν1 such that:

1. ν = ν0 + ν1 ;

2. ν0 µ (i.e. ν0 is absolutely continuous with respect to µ;)

3. ν1 ⊥ µ (i.e. ν1 and µ are singular.)

These two measures are uniquely determined.

Version: 5 Owner: Koro Author(s): Koro

396.4 Lebesgue outer measure

**Let S be some arbitrary subset of R. Let L(I) be the traditional definition of the length of
**

an interval I ⊆ R. If I = (a, b), then L(I) = b − a. Let M be the set containing

X

L(A)

A∈C

S

for any countable collection of open intervals C that covers S (that is, S ⊆ C). Then the

Lebesgue outer measure of S is defined by:

m∗ (S) = inf(M)

Note that (R, P(R), m∗ ) is “almost” a measure space. In particular:

1541

• Lebesgue outer measure is defined for any subset of R (and P(R) is a σ-algebra).

**• m∗ A > 0 for any A ⊆ R, and m∗ ∅ = 0.
**

S

• If A and B are disjoint sets, then m∗ (A B) 6 mS ∗

A + m∗P

B. More generally, if hAi i

∗

is a countable sequence of disjoint sets, then m ( Ai ) 6 m∗ Ai . This property is

known as countable subadditivity and is weaker than countable additivity. In fact,

m∗ is not countably additive.

Lebesgue outer measure has other nice properties:

• The outer measure of an interval is its length: m∗ (a, b) = b − a.

**• m∗ is translation invariant. That is, if we define A + y to be the set {x + y : x ∈ A},
**

we have m∗ A = m∗ (A + y) for any y ∈ R.

Version: 4 Owner: vampyr Author(s): vampyr

396.5 absolutely continuous

**Given two signed measures µ and ν on the same measurable space (Ω, S), we say that ν is
**

absolutely continuous with respect to µ if, for each A ∈ S such that |µ|(A) = 0, it holds

ν(A) = 0. This is usually denoted by ν µ.

Remarks.

If (ν + , ν − ) is the Jordan decomposition of ν, the following propositions are equivalent:

1. ν µ;

2. ν + µ and ν − µ;

3. |ν| kµ|.

**If ν is a finite signed measure and ν µ, the following useful property holds: for each ε > 0,
**

there is a δ > 0 such that |ν|(E) < ε whenever |µ|(E) < δ.

Version: 5 Owner: Koro Author(s): Koro

1542

396.6 counting measure

**Let (X, B) be a measurable space. We call a measure µ counting measure on X if
**

n if A has exactly n elements

µ(A ∈ B) =

∞ otherwise.

Generally, counting measure is applied on N or Z.

Version: 2 Owner: mathwizard Author(s): mathwizard, drummond

396.7 measurable set

**Let (X, F, µ) be a measure space with a sigma algebra F. A measurable set with respect
**

to µ in X is an element of F. These are also sometimes called µ-measurable sets. Any subset

Y ⊂ X with Y ∈ / F is said to be nonmeasurable with respect to µ, or non-µ-measurable.

Version: 2 Owner: mathcam Author(s): mathcam, drummond

396.8 outer regular

**Let X be a locally compact Hausdorff topological space with Borel σ–algebra B, and suppose
**

µ is a measure on (X, B). For any Borel set B ∈ B, the measure µ is said to be outer regular

on B if

µ(B) = inf {µ(U) | U ⊃ B, U open}.

We say µ is inner regular on B if

µ(B) = sup {µ(K) | K ⊂ B, K compact}.

Version: 1 Owner: djao Author(s): djao

**396.9 signed measure
**

S

A signed measure on a measurable space (Ω, S) is a function µ : S → R {+∞} which is

σ-additive and such that µ(∅) = 0.

Remarks.

1543

1. The usual (positive) measure is a particular case of signed measure, in which |µ| = µ

(see Jordan decomposition.)

2. Notice that the value −∞ is not allowed.

**3. An important example of signed measures arises from the usual measures in the follow-
**

ing way: Let (Ω, S, µ) be a measure space, and let f be a (real valued) measurable function

such that

int{x∈Ω:f (x)<0} |f |dµ < ∞.

Then a signed measure is defined by

A 7→ intA f dµ.

Version: 4 Owner: Koro Author(s): Koro

396.10 singular measure

**Two measures µ and ν in a measurableSspace (Ω, A) are called singular if there exist two
**

disjoint sets A and B in A such that A B = Ω and µ(B) = ν(A) = 0. This is denoted by

µ ⊥ ν.

Version: 4 Owner: Koro Author(s): Koro

1544

Chapter 397

**28A15 – Abstract differentiation
**

theory, differentiation of set functions

397.1 Hardy-Littlewood maximal theorem

**There is a constant K > 0 such that for each Lebesgue integrable function f ∈ L1 (Rn ), and
**

each t > 0,

K K

m({x : Mf (x) > t}) 6 kf k1 = intRn |f (x)|dx,

t t

where Mf is the Hardy-Littlewood maximal function of f .

Remark. The theorem holds for the constant K = 3n .

Version: 1 Owner: Koro Author(s): Koro

397.2 Lebesgue differentiation theorem

**Let f be a locally integrable function on Rn with Lebesgue measure m, i.e. f ∈ L1loc (Rn ).
**

Lebesgue’s differentiation theorem basically says that for almost every x, the averages

1

intQ |f (y) − f (x)|dy

m(Q)

converge to 0 when Q is a cube containing x and m(Q) → 0.

**Formally, this means that there is a set N ⊂ Rn with µ(N) = 0, such that for every x ∈
**

/N

and ε > 0, there exists δ > 0 such that, for each cube Q with x ∈ Q and m(Q) < δ, we have

1

intQ |f (y) − f (x)|dy < ε.

m(Q)

1545

For n = 1, this can be restated as an analogue of the fundamental theorem of calculus for

Lebesgue integrals. Given a x0 ∈ R,

d

intx f (t)dt = f (x)

dx x0

for almost every x.

Version: 6 Owner: Koro Author(s): Koro

397.3 Radon-Nikodym theorem

**Let µ and ν be two σ-finite measures on the same measurable space (Ω, S), such that ν µ
**

(i.e. ν is absolutely continuous with respect to µ.) Then there exists a measurable function

f , which is nonnegative and finite, such that for each A ∈ S,

ν(A) = intA f dµ.

This function is unique (any other function satisfying these conditions is equal to f µ-almost

everywhere,) and it is called the Radon-Nikodym derivative of ν with respect to µ,

dν

denoted by f = dµ .

**Remark. The theorem also holds if ν is a signed measure. Even if ν is not σ-finite the
**

theorem holds, with the exception that f is not necessarely finite.

Some properties of the Radon-Nikodym derivative

Let ν, µ, and λ be σ-finite measures in (Ω, S).

1. If ν λ and µ λ, then

d(ν + µ) dν dµ

= + µ-almost everywhere;

dλ dλ dλ

2. If ν µ λ, then

dν dν dν

= µ-almost everywhere;

dλ dµ dλ

3. If µ λ and g is a µ-integrable function, then

dµ

intΩ gdµ = intΩ g dλ;

dλ

4. If µ ν and ν µ, then −1

dµ dν

= .

dν dµ

Version: 5 Owner: Koro Author(s): Koro

1546

397.4 integral depending on a parameter

**Suppose (E, B, µ) is a measure space, suppose I is an open interval in R, and suppose we
**

are given a function

f : E × I → R,

(x, t) 7→ f (x, t),

where R is the extended real numbers. Further, suppose that for each t ∈ I, the mapping

x 7→ f (x, t) is in L1 (E). (Here, L1 (E) is the set of measurable functions f : E → R with

finite Lebesgue integral; intE |f (x)|dµ < ∞.) Then we can define a function F : I → R by

F (t) = intE f (x, t)dµ.

Continuity of F

Let t0 ∈ I. In addition to the above, suppose:

**1. For almost all x ∈ E, the mapping t 7→ f (x, t) is continuous at t = t0 .
**

2. There is a function g ∈ L1 (E) such that for almost all x ∈ E,

|f (x, t)| ≤ g(x)

for all t ∈ I.

Then F is continuous at t0 .

Differentiation under the integral sign

Suppose that the assumptions given in the introduction hold, and suppose:

**1. For almost all x ∈ E, the mapping t 7→ f (x, t) is differentiable for all t ∈ I.
**

2. There is a function g ∈ L1 (E) such that for almost all x ∈ E,

d

| f (x, t)| ≤ g(x)

dt

for all t ∈ I.

d

Then F is differentiable on I, dt

f (x, t)dµ is in L1 (E), and for all t ∈ I,

d d

F (t) = intE f (x, t)dµ. (397.4.1)

dt dt

The above results can be found in [1, 1].

1547

REFERENCES

1. F. Jones, Lebesgue Integration on Euclidean Spaces, Jones and Barlett Publishers, 1993.

2. C.D. Aliprantis, O. Burkinshaw, Principles of Real Analysis, 2nd ed., Academic Press,

1990.

Version: 1 Owner: matte Author(s): matte

1548

Chapter 398

**28A20 – Measurable and
**

nonmeasurable functions, sequences of

measurable functions, modes of

convergence

398.1 Egorov’s theorem

**Let (X, S, µ) be a measure space, and let E be a subset of X of finite measure. If fn is a
**

sequence of measurable functions converging to f almost everywhere, then for each δ > 0

there exists a set Eδ such that µ(Eδ ) < δ and fn → f uniformly on E − Eδ .

Version: 2 Owner: Koro Author(s): Koro

398.2 Fatou’s lemma

If f1 , f2 , . . . is a sequence of nonnegative measurable functions in a measure space X, then

**intX lim inf fn 6 lim inf intX fn
**

n→∞ n→∞

Version: 3 Owner: Koro Author(s): Koro

1549

398.3 Fatou-Lebesgue theorem

**Let X be a measure space. If Φ is a measurable function with intX Φ < ∞, and if f1 , f2 , . . .
**

is a sequence of measurable functions such that |fn | 6 Φ for each n, then

**g = lim inf fn and h = lim sup fn
**

n→∞ n→∞

are both integrable, and

**−∞ < intX g 6 lim inf intX fn 6 lim sup intX fn 6 intX h < ∞.
**

n→∞ k→∞

Version: 3 Owner: Koro Author(s): Koro

398.4 dominated convergence theorem

**Let X be a measure space, and let Φ, f1 , f2 , . . . be measurable functions such that intX Φ <
**

∞ and |fn | 6 Φ for each n. If fn → f almost everywhere, then f is integrable and

**lim intX fn = intX f.
**

n→∞

This theorem is a corollary of the Fatou-Lebesgue theorem.

**A possible generalization is that if {fr : r ∈ R} is a family of measurable functions such that
**

|fr | 6 |Φ| for each r ∈ R and fr −−→ f , then f is integrable and

r→0

**lim intX fr = intX f.
**

r→0

Version: 8 Owner: Koro Author(s): Koro

398.5 measurable function

**Let f : X → R̄ be a function defined on a measure space X. We say that f is measurable
**

if {x ∈ X | f (x) > a} is a measurable set for all a ∈ R.

Version: 5 Owner: vypertd Author(s): vypertd

1550

398.6 monotone convergence theorem

**Let X be a measure space, and let 0 6 f1 6 f2 6 · · · be a monotone increasing sequence
**

of nonnegative measurable functions. Let f be the function defined almost everywhere by

f (x) = limn→∞ fn (x). Then f is measurable, and

**lim intX fn = intX f.
**

n→∞

**Remark. This theorem is the first of several theorems which allow us to “exchange inte-
**

gration and limits”. It requires the use of the Lebesgue integral: with the Riemann integral,

we cannot even formulate the theorem, lacking, as we do, the concept of “almost every-

where”. For instance, the characteristic function of the rational numbers in [0, 1] is not

Riemann integrable, despite being the limit of an increasing sequence of Riemann integrable

functions.

Version: 5 Owner: Koro Author(s): Koro, ariels

398.7 proof of Egorov’s theorem

**Let Ei,j = {x ∈ E : |fj (x) − f (x)| < 1/i}. Since fn → f almost everywhere, there is a set
**

S with µ(S) = 0 such that, given i ∈ N and x ∈ E − S, there is m ∈ N such that j > m

implies |fj (x) − f (x)| < 1/i. This can be expressed by

[ \

E−S ⊂ Ei,j ,

m∈N j>m

**or, in other words, \ [
**

(E − Ei,j ) ⊂ S.

m∈N j>m

S

Since { j>m (E − Ei,j )}m∈N is a decreasing nested sequence of sets, each of which has finite

measure, and such that its intersection has measure 0, by continuity from above we know

that [

µ( (E − Ei,j )) −−−→ 0.

m→∞

j>m

**Therefore, for each i ∈ N, we can choose mi such that
**

[ δ

µ( (E − Ei,j )) < .

j>mi

2i

Let [ [

Eδ = (E − Ei,j ).

i∈N j>mi

1551

Then ∞ ∞

X [ X δ

µ(Eδ ) 6 µ( (E − Ei,j )) < = δ.

i=1 j>m i=1

2i

i

**We claim that fn → f uniformly on E −Eδ . In fact, given ε > 0, choose n such that 1/n < ε.
**

If x ∈ E − Eδ , we have \ \

x∈ Ei,j ,

i∈N j>mi

**which in particular implies that, if j > mn , x ∈ En,j ; that is, |fj (x) − f (x)| < 1/n < ε.
**

Hence, for each xε > 0 there is N (which is given by mn above) such that j > N implies

|fj (x) − f (x)| < ε for each x ∈ E − Eδ , as required. This completes the proof.

Version: 3 Owner: Koro Author(s): Koro

398.8 proof of Fatou’s lemma

Let f (x) = lim inf n→∞ fn (x) and let gn (x) = inf k≥n fk (x) so that we have

f (x) = sup gn (x).

n

**As gn is an increasing sequence of measurable nonnegative functions we can apply the
**

monotone convergence theorem to obtain

**intX f dµ = lim intX gn dµ.
**

n→∞

On the other hand, being gn ≤ fn , we conclude by observing

**lim intX gn dµ = lim inf intX gn dµ ≤ lim inf intX fn dµ.
**

n→∞ n→∞ n→∞

Version: 1 Owner: paolini Author(s): paolini

398.9 proof of Fatou-Lebesgue theorem

**By Fatou’s lemma we have
**

intX g ≤ lim inf intX fn

n→∞

and (recall that lim sup f = − lim inf −f )

**lim sup intX fn ≤ intX h.
**

n→∞

1552

On the other hand by the properties of lim inf and lim sup we have

g ≥ −Φ, f ≤Φ

and hence

intX g ≥ intX − Φ > −∞, intX h ≤ intX Φ < +∞.

Version: 1 Owner: paolini Author(s): paolini

398.10 proof of dominated convergence theorem

It is not difficult to prove that f is measurable. In fact we can write

**f (x) = sup inf k≥n fk (x)
**

n

and we know that measurable functions are closed under the sup and inf operation.

**Consider the sequence gn (x) = 2Φ(x) − |f (x) − fn (x)|. clearly gn are nonnegative functions
**

since f − fn ≤ 2Φ. So, applying Fatou’s lemma, we obtain

**lim intX |f − fn | dµ ≤ lim sup intX |f − fn | dµ
**

n→∞ n→∞

= − lim inf intX − |f − fn | dµ

n→∞

= intX 2Φ dµ − lim inf intX 2Φ − |f − fn | dµ

n→∞

≤ intX 2Φ dµ − intX 2Φ − lim sup |f − fn | dµ

n→∞

= intX 2Φ dµ − intX 2Φ dµ = 0.

Version: 1 Owner: paolini Author(s): paolini

398.11 proof of monotone convergence theorem

**It is enough to prove the following
**

S

Theorem 7. Let (X, µ) be a measurable space and let fk : X → R {+∞} be a monotone increasing

sequence of positive measurable functions (i.e. 0 ≤ f1 ≤ f2 ≤ . . .). Then f (x) = limk→∞ fk (x)

is measurable and

lim intX fk dµ = intX f (x).

n→∞

1553

First of all by the monotonicity of the sequence we have

f (x) = sup fk (x)

k

**hence we know that f is measurable. Moreover being fk ≤ f for all k, by the monotonicity
**

of the integral, we immediately get

**sup intX fk dµ ≤ intX f (x) dµ.
**

k

So take any simple measurable function s such that 0 ≤ s ≤ f . Given also α < 1 define

Ek = {x ∈ X : fk (x) ≥ αs(x)}.

**The sequence Ek is an increasing sequence of measurable sets. Moreover the union of all Ek
**

is the whole space X since limk→∞ fk (x) = f (x) ≥ s(x) > αs(x). Moreover it holds

intX fk dµ ≥ intEk fk dµ ≥ αintEk s dµ.

**Being s a simple measurable function it is easy to check that E 7→ intE s dµ is a measure and
**

hence

sup intX fk dµ ≥ αintX s dµ.

k

But this last inequality holds for every α < 1 and for all simple measurable functions s with

s ≤ f . Hence by the definition of Lebesgue integral

**sup intk fk dµ ≥ intX f dµ
**

k

which completes the proof.

Version: 1 Owner: paolini Author(s): paolini

1554

Chapter 399

**28A25 – Integration with respect to
**

measures and other set functions

399.1 L∞(X, dµ)

**The L∞ space, L∞ (X, dµ), is a vector space consisting of equivalence classes of functions
**

f : X → C with norm given by

kf k∞ = ess sup |f (t)| ,

the essential supremum of |f |. Additionally, we require that kf k∞ < ∞.

**The equivalence classes of L∞ (X, dµ) are given by saying that f, g : X → C are equivalent
**

iff f and g differ on a set of µ measure zero.

Version: 3 Owner: ack Author(s): bbukh, ack, apmxi

399.2 Hardy-Littlewood maximal operator

**The Hardy-Littlewood maximal operator in Rn is an operator defined on L1loc (Rn ) (the
**

space of locally integrable functions in Rn with the Lebesgue measure) which maps each

locally integrable function f to another function Mf , defined for each x ∈ Rn by

1

Mf (x) = sup intQ |f (y)|dy,

Q m(Q)

**where the supremum is taken over all cubes Q containing x. This function is lower semicontinuous
**

(and hence measurable), and it is called the Hardy-Littlewood maximal function of f .

1555

The operator M is sublinear, which means that

M(af + bg) 6 |a|Mf + |b|Mg

for each pair of locally integrable functions f, g and scalars a, b.

Version: 3 Owner: Koro Author(s): Koro

**399.3 Lebesgue integral
**

S

The integral of a measurable function f : X → R {±∞} on a measure space (X, B, µ) is

written

intX f dµ or just intf. (399.3.1)

It is defined via the following steps:

**• If f = f rm[o]−−A is the characteristic function of a set A ∈ B, then set
**

intX f rm[o]−−A dµ := µ(A). (399.3.2)

**• If f is a simple function (i.e. if f can be written as
**

n

X

f= ck f rm[o]−−Ak , ck ∈ R (399.3.3)

k=1

**for some finite collection Ak ∈ B), then define
**

n

X n

X

intX f dµ := ck intX f rm[o]−−Ak dµ = ck µ(Ak ). (399.3.4)

k=1 k=1

**• If f is a nonnegative measurable function (possibly attaining the value ∞ at some
**

points), then we define

intX f dµ := sup {intX h dµ : h is simple and h(x) ≤ f (x) for all x ∈ X} . (399.3.5)

**• For any measurable function f (possibly attaining the values ∞ or −∞ at some points),
**

write f = f + − f − where

f + := max(f, 0) and f − := max(−f, 0), (399.3.6)

and define the integral of f as

intX f dµ := intX f + dµ − intX f − dµ, (399.3.7)

provided that intX f + dµ and intX f − dµ are not both ∞.

1556

If µ is Lebesgue measure and X is any interval in Rn then the integral is called the Lebesgue

integral. If the Lebesgue integral of a function f on a set A exists, f is said to be Lebesgue

integrable. The Lebesgue integral equals the Riemann integral everywhere the latter is

defined; the advantage to the Lebesgue integral is that many Lebesgue-integrable functions

are not Riemann-integrable. For example, the Riemann integral of the characteristic function

of the rationals in [0, 1] is undefined, while the Lebesgue integral of this function is simply

the measure of the rationals in [0, 1], which is 0.

Version: 12 Owner: djao Author(s): djao, drummond

1557

Chapter 400

**28A60 – Measures on Boolean rings,
**

measure algebras

400.1 σ-algebra

Let X be a set. A σ-algebra is a collection M of subsets of X such that

• X∈M

• If A ∈ M then X − A ∈ M.

**• If A1 , A2 , A3 , . . . is a countable subcollection of M, that is, Aj ∈ M for j = 1, 2, 3, . . .
**

(the subcollection can be finite) then the union of all of them is also in M:

∞

[

Ai ∈ M.

j=1

Version: 3 Owner: drini Author(s): drini, apmxi

400.2 σ-algebra

**Given a set E, a sigma algebra (or σ–algebra) in E is a collection B(E) of subsets of E such
**

that:

• ∅ ∈ B(E)

• Any countable union of elements of B(E) is in B(E)

1558

• The complement of any element of B(E) in E is in B(E)

**Given any collection C of subsets of B(E), the σ–algebra generated by C is defined to be
**

the smallest σ–algebra in E containing C.

Version: 5 Owner: djao Author(s): djao

400.3 algebra

Given a set E, an algebra in E is a collection B(E) of subsets of E such that:

• ∅ ∈ B(E)

• Any finite union of elements of B(E) is in B(E)

• The complement of any element of B(E) in E is in B(E)

**Given any collection C of subsets of B(E), the algebra generated by C is defined to be the
**

smallest algebra in E containing C.

Version: 2 Owner: djao Author(s): djao

400.4 measurable set (for outer measure)

**Definition [1, 2, 1] Let µ∗ be an outer measure on a set X. A set E ⊂ X is said to be
**

measurable, or µ∗ -measurable, if for all A ⊂ X, we have

\ \

µ∗ (A) = µ∗ (A E) + µ∗ (A E {). (400.4.1)

**Remark If A, E ⊂ X, we have, from the properties of the outer measure,
**

\ [

µ∗ (A) = µ∗ A (E E {)

\ [ \

= µ∗ (A E) (A E {)

\ \

≤ µ∗ (A E) + µ∗ (A E {)).

**Hence equation (399.4.1) is equivalent to the inequality [1, 2, 1]
**

\ \

µ∗ (A) ≥ µ(A E) + µ(A E {).

1559

Of course, this inequality is trivially satisfied if µ∗ (A) = ∞. Thus a set E ⊂ X is µ-

measurable in X if and only if the above inequality holds for all A ⊂ X for which µ∗ (A) < ∞

[1].

**Theorem [Carathéodory’s theorem] [1, 2, 1] Suppose µ∗ is an outer measure on a set
**

X, and suppose M is the set of all µ∗ -measurable sets in X. Then M is a σ-algebra, and µ∗

restricted to M is a measure (on M).

Example Let µ∗ be an outer measure on a set X.

∗ ∗

1. Any null set (a set E with

T µ (E) = 0) is measurable.

T Indeed, suppose µT (E) = 0, and

A ⊂ X. Then, sinceTA E ⊂ E, we have µ (A E) = 0, and since A E { ⊂ A, we

∗

have µ∗ (A) ≥ µ∗ (A E {), so

\

µ∗ (A) ≥ µ∗ (A E)

\ \

= µ∗ (A E) + µ∗ (A E {).

Thus E is measurable.

S∞

2. If {Bi }∞

i=1 is a countable collection of null sets, then i=1 Bi is a null set. This follows

directly from the last property of the outer measure.

REFERENCES

1. A. Mukherjea, K. Pothoven, Real and Functional analysis, Plenum press, 1978.

2. A. Friedman, Foundations of Modern Analysis, Dover publications, 1982.

3. G.B. Folland, Real Analysis: Modern Techniques and Their Applications, 2nd ed,

John Wiley & Sons, Inc., 1999.

Version: 1 Owner: matte Author(s): matte

1560

Chapter 401

**28A75 – Length, area, volume, other
**

geometric measure theory

401.1 Lebesgue density theorem

**Let µ be the Lebesgue measure on R. If µ(Y ) > 0, then there exists X ⊂ Y such that
**

µ(Y − X) = 0 and for all x ∈ X

T

µ(X [x − , x + ])

lim = 1.

→+0 2

Version: 2 Owner: bbukh Author(s): bbukh

1561

Chapter 402

28A80 – Fractals

402.1 Cantor set

**The Cantor set C is the canonical example of an uncountable set of measure zero. We
**

construct C as follows.

**Begin with the unit interval C0 = [0, 1], and remove the open segment R1 := ( 31 , 23 ) from the
**

middle. We define C1 as the two remaining pieces

1 [ 2

C1 := C0 R1 = 0, ,0 (402.1.1)

3 3

Now repeat the process on each remaining segment, removing the open set

1 2 [ 7 8

R2 := , , (402.1.2)

9 9 9 9

to form the four-piece set

1 [ 2 1 [ 2 7 [ 8

C2 := C1 R2 = 0, , , ,1 (402.1.3)

9 9 3 3 9 16

Continue the process, forming C3 , C4 , . . . Note that Ck has 2k pieces.

Figure 402.1: The sets C0 through C5 in the construction of the Cantor set

Also note that at each step, the endpoints of each closed segment will stay in the set forever—

e.g., the point 32 isn’t touched as we remove sets.

1562

The Cantor set is defined as

∞

\ ∞

[

C := Ck = C0 \ Rn (402.1.4)

k=1 n=1

**Cardinality of the Cantor set To establish cardinality, we want a bijection between
**

some set whose cardinality we know (e.g. Z, R) and the points in the Cantor set. We’ll be

aggressive and try the reals.

Start at C1 , which has two pieces. Mark the left-hand segment “0” and the right-hand

segment “1”. Then continue to C2 , and consider only the leftmost pair. Again, mark the

segments “0” and “1”, and do the same for the rightmost pair.

Keep doing this all the way down the Ck , starting at the left side and marking the segments

0, 1, 0, 1, 0, 1 as you encounter them, until you’ve labeled the entire Cantor set.

Now, pick a path through the tree starting at C0 and going left-left-right-left. . . and so on.

Mark a decimal point for C0 , and record the zeros and ones as you proceed. Each path has a

unique number based on your decision at each step. For example, the figure represents your

choice of left-left-right-left-right at the first five steps, representing the number beginning

0.00101...

Every point in the Cantor set will have a unique address dependent solely on the pattern

Figure 402.2: One possible path through C5 : 0.00101

of lefts and rights, 0’s and 1’s, required to reach it. Each point thus has a unique number,

the real number whose binary expansion is that sequence of zeros and ones. Every infinite

stream of binary digits can be found among these paths, and in fact the binary expansion of

every real number is a path to a unique point in the Cantor set.

**Some caution is justified, as two binary expansions may refer to the same real number; for
**

example, 0.011111 . . . = 0.100000 . . . = 21 . However, each one of these duplicates must corre-

spond to a rational number. To see this, suppose we have a number x in [0, 1] whose binary

expansion becomes all zeros or all ones at digit k (both are the same number, remember).

Then we can multiply that number by 2k and get 1, so it must be a (binary) rational number.

There are only countably many rationals, and not even all of those are the double-covered

numbers we’re worried about (see, e.g., 13 = 0.0101010 . . .), so we have at most countably

many duplicated reals. Thus, the cardinality of the 0.Cantor set is equal to that of the reals.

(If we want to be really picky, map (0, 1) to the reals with, say, f (x) = 1/x + 1/(x − 1), and

the end points really don’t matter much.)

**Return, for a moment, to the earlier observation that numbers such as 31 , 29 , the endpoints
**

of deleted intervals, are themselves never deleted. In particluar, consider the first deleted

1563

interval: the ternary expansions of its constituent numbers are precisely those that begin

0.1, and proceed thence with at least one non-zero “ternary” digit (just digit for us) further

along. Note also that the point 13 , with ternary expansion 0.1, may also be written 0.02̇

(or 0.02̄), which has no digits 1. Similar descriptions apply to further deleted intervals.

The result is that the cantor set is precisely those numbers in the set [0, 1] whose ternary

expansion contains no digits 1.

**Measure of the Cantor set Let µ be Lebesgue measure. The measure of the sets Rk
**

that we remove during the construction of the Cantor set are

2 1 1

µ(R1 ) = − = (402.1.5)

3 3 3

2 1 8 7 2

µ(R2 ) = − + − = (402.1.6)

9 9 9 9 9

..

. (402.1.7)

k

X 2n−1

µ(Rk ) = (402.1.8)

n=1

3n

Note that the R’s are disjoint, which will allow us to sum their measures without worry. In

the limit k → ∞, this gives us

∞

! ∞

[ X 2n−1

µ Rn = = 1. (402.1.9)

n=1 n=1

3n

**But we have µ(C0 ) = 1 as well, so this means
**

∞

! ∞

[ X 1

µ(C) = µ C0 \ Rn = µ(C0) − = 1 − 1 = 0. (402.1.10)

n=1 n=1

2n

Thus we have seen that the measure of C is zero (though see below for more on this topic).

How many points are there in C? Lots, as we shall see.

So we have a set of measure zero (very tiny) with uncountably many points (very big). This

non-intuitive result is what makes Cantor sets so interesting.

**Cantor sets with positive measure clearly, Cantor sets can be constructed for all sorts
**

of “removals”—we can remove middle halves, or thirds, or any amount 1r , r > 1 we like. All

of these Cantor sets have measure zero, since at each step n we end up with

n

1

Ln = 1 − (402.1.11)

r

1564

of what we started with, and limn→∞ Ln = 0 for any r > 1. With apologies, the figure above

is drawn for the case r = 2, rather than the r = 3 which seems to be the publically favored

example.

**However, it is possible to construct Cantor sets with positive measure as well; the key is to
**

remove less and less as we proceed. These Cantor sets have the same “shape” (topology) as

the Cantor set we first constructed, and the same cardinality, but a different “size.”

**Again, start with the unit interval for C0 , and choose a number 0 < p < 1. Let
**

2−p 2+p

R1 := , (402.1.12)

4 4

**which has length (measure) p2 . Again, define C1 := C0 \ R1 . Now define
**

2 − p 2 + p [ 14 − p 14 + p

R2 := , , (402.1.13)

16 16 16 16

**which has measure p4 . Continue as before, such that each Rk has measure 2pk ; note again
**

that all the Rk are disjoint. The resulting Cantor set has measure

∞

! ∞ ∞

[ X X

µ C0 \ Rn = 1 − µ(Rn ) = 1 − p 2−n = 1 − p > 0.

n=1 n=1 n=1

**Thus we have a whole family of Cantor sets of positive measure to accompany their vanishing
**

brethren.

Version: 19 Owner: drini Author(s): drini, quincynoodles, drummond

402.2 Hausdorff dimension

**Let Θ be a bounded subset of Rn let NΘ () be the minimum number of balls of radius
**

required to cover Θ. Then define the Hausdorff dimension dH of Θ to be

log NΘ ()

dH (Θ) := − lim .

→0 log

**Hausdorff dimension is easy to calculate for simple objects like the Sierpinski gasket or a
**

Koch curve. Each of these may be covered with a collection of scaled-down copies of itself.

In fact, in the case of the Sierpinski gasket, one can take the individual triangles in each

approximation as balls in the covering. At stage n, there are 3n triangles of radius 21n , and

so the Hausdorff dimension of the Sierpinski triangle is − nnlog

log 3

1/2

log 3

= log 2

.

1565

From some notes from Koro This definition can be extended to a general metric space

X with distance function d.

**Define the diameter |C| of a bounded subset C of X to be supx,y∈C d(x, y), and define a
**

countable r-cover S

of X to be a collection of subsets Ci of X indexed by some countable set

I, such that X = i∈I Ci . We also define the handy function

X

HrD (X) = inf |Ci |D

i∈I

**where the infimum is over all countable r-covers of X. The Hausdorf dimension of X may
**

then be defined as

dH (X) = inf{D | lim HrD (X) = 0}.

r→0

n

When X is a subset of R with any restricted norm-induced metric, then this definition

reduces to that given above.

Version: 8 Owner: drini Author(s): drini, quincynoodles

402.3 Koch curve

**A Koch curve is a fractal generated by a replacement rule. This rule is, at each step, to
**

replace the middle 1/3 of each line segment with two sides of a right triangle having sides of

length equal to the replaced segment. Two applications of this rule on a single line segment

gives us:

**To generate the Koch curve, the rule is applied indefinitely, with a starting line segment.
**

Note that, if the length of the initial line segment is l, the length LK of the Koch curve at

the nth step will be

n

4

LK = l

3

**This quantity increases without bound; hence the Koch curve has infinite length. However,
**

the curve still bounds a finite area. We can prove this by noting that in each step, we add an

amount of area equal to the area of all the equilateral triangles we have just created. We can

bound the area of each triangle of side length s by s2 (the square containing the triangle.)

Hence, at step n, the area AK ”under” the Koch curve (assuming l = 1) is

1566

Figure 402.3: Sierpinski gasket stage 0, a single triangle

Figure 402.4: Stage 1, three triangles

2 2 2

1 1 1

AK < +3 +9 +···

3 9 27

Xn

1

= i−1

i=1

3

**but this is a geometric series of ratio less than one, so it converges. Hence a Koch curve has
**

infinite length and bounds a finite area.

**A Koch snowflake is the figure generated by applying the Koch replacement rule to an
**

equilateral triangle indefinitely.

Version: 3 Owner: akrowne Author(s): akrowne

402.4 Sierpinski gasket

**Let S0 be a triangular area, and define Sn+1 to be obtained from Sn by replacing each trian-
**

gular area in Sn with three similar and similarly oriented triangular areas each intersecting

with each of the other two at exactly one vertex, each one half the linear scale of the orrig-

inal in size. The limiting set as n → ∞ (alternately the intersection of all these sets) is a

Sierpinski gasket, also known as a Sierpinski triangle.

Version: 3 Owner: quincynoodles Author(s): quincynoodles

402.5 fractal

**Option 1: Some equvialence class of subsets of Rn . A usual equivalence is postulated when
**

some generalised ”distance” is zero. For example, let F, G ⊂ Rn , and let d(x, y) be the usual

distance (x, y ∈ R). Define the distance D between F and G as

**D(F, G) := inf f ∈F sup d(f, g) + inf g∈G sup d(f, g)
**

g∈G f ∈F

Figure 402.5: Stage 2, nine triangles

1567

Figure 402.6: Stage n, 3n triangles

Then in this case we have, as fractals, that Q and R are equivalent.

**Option 2: A subset of Rn with non-integral Hausdorff dimension. Examples: (we think) the
**

coast of Britain, a Koch snowflake.

**Option 3: A “self-similar object”. That is, one which can be covered by copies of itself using
**

a set of (usually two or more) transformation mappings. Another way to say this would

be “an object with a discrete approximate scaling symmetry.” Example: A square region,

a Koch curve, a fern frond. This isn’t much different from Option 1 because of the collage

theorem.

**A cursory description of some relationships between options 2 and 3 is given towards the
**

end of the entry on Hausdorff dimension. The use of option 1 is that it permits one to talk

about how ”close” two fractals are to one another. This becomes quite handy when one

wants to talk about approximating fractals, especially approximating option 3 type fractals

with pictures that can be drawn in finite time. A simple example: one can talk about how

close one of the line drawings in the Koch curve entry is to an actual Koch curve.

Version: 7 Owner: quincynoodles Author(s): quincynoodles

1568

Chapter 403

28Axx – Classical measure theory

403.1 Vitali’s Theorem

There exists a set V ⊂ [0, 1] which is not Lebesgue measurable

Version: 1 Owner: paolini Author(s): paolini

403.2 proof of Vitali’s Theorem

**Consider the equivalence relation in [0, 1) given by
**

x∼y ⇔ x−y ∈Q

and let F be the family of all equivalence classes of ∼. Let V be a section of F i.e. put in V

an element for each equivalence class of ∼ (notice that we are using the axiom of choice).

T

Given q ∈ Q [0, 1) define

\ [ \

Vq = ((V + q) [0, 1)) ((V + q − 1) [0, 1))

that is Vq is obtained translating V by a quantity q to the right and then cutting the piece

which goes beyond the point 1 and putting it on the left, starting from 0.

**Now notice that given x ∈ [0, 1) there
**

T exists y ∈ V such that x ∼ y (because V is a section

of ∼) and hence there exists q ∈ Q [0, 1) such that x ∈ Vq . So

[

Vq = [0, 1).

T

q∈Q [0,1)

T

Moreover all the Vq are disjoint. In fact if x ∈ Vq Vp then x − q (modulus [0, 1)) and x − p

are both in V which is not possible since they differ by a rational quantity q −p (or q −p + 1).

1569

Now if V is Lebesgue measurable, clearly also Vq are measurable and µ(Vq ) = µ(V ). Moreover

by the countable additivity of µ we have

X X

µ([0, 1)) = µ(Vq ) = µ(V ).

q

T

q∈Q [0,1)

So if µ(V ) = 0 we had µ([0, 1)) = 0 and if µ(V ) > 0 we had µ([0, 1)) = +∞.

So the only possibility is that V is not Lebesgue measurable.

Version: 1 Owner: paolini Author(s): paolini

1570

Chapter 404

**28B15 – Set functions, measures and
**

integrals with values in ordered spaces

404.1 Lp-space

**Definition Let (X, B, µ) be a measure space. The Lp -norm of a function f : X → R is
**

defined as 1

kf kp := (intX |f |p dµ) p (404.1.1)

when the integral exists. The set of functions with finite Lp -norm form a vector space V

with the usual pointwise addition and scalar multiplication of functions. In particular, the

set of functions with zero Lp -norm form a linear subspace of V , which for this article will be

called K. We are then interested in the quotient space V /K, which consists of real functions

on X with finite Lp -norm, identified up to equivalence almost everywhere. This quotient

space is the real Lp -space on X.

Theorem The vector space V /K is complete with respect to the Lp norm.

**The space L∞ . The space L∞ is somewhat special, and may be defined without explicit
**

reference to an integral. First, the L∞ -norm of f is defined to be the essential supremum of

|f |:

kf k∞ := ess sup |f | = inf {a ∈ R : µ({x : |f (x)| > a}) = 0} (404.1.2)

The definitions of V , K, and L∞ then proceed as above. Functions in L∞ are also called

essentially bounded.

**Example Let X = [0, 1] and f (x) = √1 . Then f ∈ L1 (X) but f ∈
**

/ L2 (X).

x

Version: 18 Owner: mathcam Author(s): Manoj, quincynoodles, drummond

1571

404.2 locally integrable function

**Definition [4, 1, 2] Suppose that U is an open set in Rn , and f : U → C is a Lebesgue integrable
**

function. If the Lebesgue integral

intK |f |dx

is finite for all compact subsets K in U, then f is locally integrable. The set of all such

functions is denoted by L1loc (U).

Example

1. L1 (U) ⊂ L1loc (U), where L1 (U) is the set of integrable functions.

**Theorem Suppose f and g are locally integrable functions on an open subset U ⊂ Rn , and
**

suppose that

intU f φdx = intU gφdx

for all smooth functions with compact support φ ∈ C0∞ (U). Then f = g almost everywhere.

**A proof based on the Lebesgue differentiation theorem is given in [4] pp. 15. Another proof
**

is given in [2] pp. 276.

REFERENCES

1. L. Hörmander, The Analysis of Linear Partial Differential Operators I, (Distribution

theory and Fourier Analysis), 2nd ed, Springer-Verlag, 1990.

2. G.B. Folland, Real Analysis: Modern Techniques and Their Applications, 2nd ed,

John Wiley & Sons, Inc., 1999.

3. S. Lang, Analysis II, Addison-Wesley Publishing Company Inc., 1969.

Version: 3 Owner: matte Author(s): matte

1572

Chapter 405

**28C05 – Integration theory via linear
**

functionals (Radon measures, Daniell

integrals, etc.), representing set

functions and measures

405.1 Haar integral

**Let Γ be a locally compact topological group and C be the algebra of all continuous real-
**

valued functions on Γ with compact support. In addition we define C+ to be the set of

non-negative functions that belong to C. The Haar integral is a real linear map I of C into

the field of the real number for Γ if it satisfies:

• I is not the zero map

• I only takes non-negative values on C+

**• I has the following property I(γ · f ) = I(f ) for all elements f of C and all element γ
**

of Γ.

The Haar integral may be denoted in the following way (there are also other ways):

intγ∈Γ f (γ) or intΓ f or intΓ f dγ or I(f )

**In order for the Haar intergral to exists and to be unique, the following conditions are
**

necessary and sufficient: That there exists a real-values function I + on C+ satisfying the

following condition:

1573

1. (Linearity).I + (λf + µg) = λI + (f ) + µI + (g) where f, g ∈ C+ and λ, µ ∈ R+ .

2. (Positivity). If f (γ) > 0 for all γ ∈ Γ then I + (f (γ)) > 0.

3. (Translation-Invariance). I(f (δγ)) = I(f (γ)) for any fixed δ ∈ Γ and every f in C+ .

**An additional property is if Γ is a compact group then the Haar integral has right translation-
**

invariance: intγ∈Γ f (γδ) = intγ∈Γ f (γ) for any fixed δ ∈ Γ. In addition we can define nor-

malized Haar integral to be intΓ 1 = 1 since Γ is compact, it implies that intΓ 1 is finite.

(The proof for existence and uniqueness of the Haar integral is presented in [PV] on page

9.)

( the information of this entry is in part quoted and paraphrased from [GSS])

REFERENCES

[GSS] Golubsitsky, Matin. Stewart, Ian. Schaeffer, G. David.: Singularities and Groups in Bifur-

cation Theory (Volume II). Springer-Verlag, New York, 1988.

[HG] Gochschild, G.: The Structure of Lie Groups. Holden-Day, San Francisco, 1965.

Version: 4 Owner: Daume Author(s): Daume

1574

Chapter 406

**28C10 – Set functions and measures
**

on topological groups, Haar measures,

invariant measures

406.1 Haar measure

406.1.1 Definition of Haar measures

**Let G be a locally compact topological group. A left Haar measure on G is a measure µ on
**

the Borel sigma algebra B of G which is:

1. outer regular on all Borel sets B ∈ B

2. inner regular on all open sets U ⊂ G

3. finite on all compact sets K ⊂ G

4. invariant under left translation: µ(gB) = µ(B) for all Borel sets B ∈ B

**A right Haar measure on G is defined similarly, except with left translation invariance re-
**

placed by right translation invariance (µ(Bg) = µ(B) for all Borel sets B ∈ B). A bi–

invariant Haar measure is a Haar measure that is both left invariant and right invariant.

1575

406.1.2 Existence of Haar measures

**For any finite group G, the counting measure on G is a bi–invariant Haar measure. More
**

generally, every locally compact topological group G has a left 1 Haar measure µ, which is

unique up to scalar multiples. The Haar measure plays an important role in the development

of Fourier analysis and representation theory on locally compact groups such as Lie groups

and profinite groups.

Version: 1 Owner: djao Author(s): djao

1

G also has a right Haar measure, although the right and left Haar measures on G are not necessarily

equal unless G is abelian.

1576

Chapter 407

**28C20 – Set functions and measures
**

and integrals in infinite-dimensional

spaces (Wiener measure, Gaussian

measure, etc.)

407.1 essential supremum

**Let (X, B, µ) be a measure space and let f : X → R be a function. The essential supre-
**

mum of f is the smallest number a ∈ R for which f only exceeds a on a set of measure zero.

This allows us to generalize the maximum of a function in a useful way.

More formally, we define ess sup f as follows. Let a ∈ R, and define

Ma = {x : f (x) > a} , (407.1.1)

the subset of X where f (x) is greater than a. Then let

A0 = {a ∈ R : µ(Ma ) = 0} , (407.1.2)

**the set of real numbers for which Ma has measure zero. If A0 = ∅, then the essential
**

supremum is defined to be ∞. Otherwise, the essential supremum of f is

ess sup f := infA0 . (407.1.3)

Version: 1 Owner: drummond Author(s): drummond

1577

Chapter 408

**28D05 – Measure-preserving
**

transformations

408.1 measure-preserving

**Let (X, B, µ) be a measure space, and T : X → X be a (possibly non-invertible) measurable
**

transformation. We call T measure-preserving if for all A ∈ B,

µ(T −1(A)) = µ(A),

where T −1 (A) is defined to be the set of points x ∈ X such that T (x) ∈ A.

A measure-preserving transformation is also called an endomorphism of the measure space.

Version: 5 Owner: mathcam Author(s): mathcam, drummond

1578

Chapter 409

**30-00 – General reference works
**

(handbooks, dictionaries,

bibliographies, etc.)

409.1 domain

A non-empty open set in C is called a domain.

**The topology considered is the Euclidean one (viewing C as R2 ). So we have that for a
**

domain D being connected is equivalent to being path-connected.

**Since we have that every component of a domain D will be a region, we have that every
**

domain has at most countably many components.

Version: 4 Owner: drini Author(s): drini

409.2 region

A region is a connected domain.

**Since every domain of C can be seen as the union of countably many components and each
**

component is a region, we have that regions play a major role in complex analysis.

Version: 2 Owner: drini Author(s): drini

1579

409.3 regular region

**Let E be a n-dimensional Euclidean space with the topology induced by the Euclidean metric.
**

Then a set in E is a regular region, if it can be written as the closure of a non-empty region

with a piecewise smooth boundary.

Version: 10 Owner: ottocolori Author(s): ottocolori

409.4 topology of the complex plane

**The usual topology for the complex plane C is the topology induced by the metric d(x, y) =
**

|x − y| for x, y ∈ C. Here, | · | is the complex modulus.

**If we identify R2 and C, it is clear that the above topology coincides with topology induced
**

by the Euclidean metric on R2 .

Version: 1 Owner: matte Author(s): matte

1580

Chapter 410

**30-XX – Functions of a complex
**

variable

410.1 z0 is a pole of f

Let f be an analytic function on a punctured neighborhood of x0 ∈ C, that is, f analytic on

{z ∈ C : 0 < |z − x0 | < ε}

**for some ε > 0 and such that
**

lim f = ∞.

z→z0

We say then that x0 is a pole for f .

Version: 2 Owner: drini Author(s): drini, apmxi

1581

Chapter 411

30A99 – Miscellaneous

411.1 Riemann mapping theorem

**Let U be a simply connected open proper subset of C, and let a ∈ U. There is a unique
**

analytic function f : U → C such that

1. f (a) = 0, and f 0 (a) is real and positive;

2. f is injective;

3. f (U) = {z ∈ C : |z| < 1}.

**Remark. As a consequence of this theorem, any two simply connected regions, none of
**

which is the whole plane, are conformally equivalent.

Version: 2 Owner: Koro Author(s): Koro

**411.2 Runge’s theorem
**

S

Let K be a compact subset of C, and let E be a subset of C∞ = C {∞} (the extended

complex plane) which intersects every connected component of C∞ − K. If f is an analytic

function in an open set containing K, given ε > 0, there is a rational function R(z) whose

only poles are in E, such that |f (z) − R(z)| < ε for all z ∈ K.

Version: 2 Owner: Koro Author(s): Koro

1582

411.3 Weierstrass M-test

**Let X be a topological space, {fn }n∈N a sequence of real or complex valued functions on
**

X and {Mn }n∈N a sequence of non-negative real Pnumbers. Suppose that, for eachP∞ n∈ N

∞

and x ∈ X, we have |fn (x)| ≤ Mn . Then f = n=1 fn converges uniformly if n=1 Mn

converges.

Version: 8 Owner: vypertd Author(s): vypertd, igor

411.4 annulus

Briefly, an annulus is the region bounded between two (usually concentric) circles.

**An open annulus, or just annulus for short, is a domain in the complex plane of the form
**

A = Aw (r, R) = {z ∈ C | r < |z − w| < R},

where w is an abitrary complex number, and r and R are real numbers with 0 < r < R.

Such a set is often called an annular region.

**More generally, one can allow r = 0 or R = ∞. (This makes sense for the purposes of the
**

bound on |z − w| above.) This would make an annulus include the cases of a punctured disc,

and some unbounded domains.

**Analogously, one can define a closed annulus to be a set of the form
**

A = Aw (r, R) = {z ∈ C | r 6 |z − w| 6 R},

where w ∈ C, and r and R are real numbers with 0 < r < R.

**One can show that two annuli Dw (r, R) and Dw0 (r 0 , R0 ) are conformally equivalent if and
**

only if R/r = R0 /r 0 . More generally, the complement of any closed disk in an open disk is

conformally equivalen to precisely one annulus of the form D0 (r, 1).

Version: 1 Owner: jay Author(s): jay

411.5 conformally equivalent

**A region G is conformally equivalent to a set S if there is an analytic bijective function
**

mapping G to S.

Conformal equivalence is an equivalence relation.

Version: 1 Owner: Koro Author(s): Koro

1583

411.6 contour integral

**Let f be a complex-valued function defined on the image of a curve α: [a, b] → C, let
**

P = {a0 , ..., an } be a partition of [a, b]. If the sum

n

X

f (zi )(α(ai) − α(ai−1 ))

i=1

**where zi is some point α(ti ) such that ai−1 6 ti 6 ai , tends to a unique limit l as n tends
**

to infinity and the greatest of the numbers ai − ai−1 tends to zero, then we say that the

contour integral of f along α exists and has value l. The contour integral is denoted by

intα f (z)dz

Note

**(i) If Im(α) is a segment of the real axis, then this definition reduces to that of the
**

Riemann integral of f(x) between α(a) and α(b)

**(ii) An alternative definition, making use of the Riemann-Stieltjes integral, is based on the
**

fact that the definition of this can be extended without any other changes in the wording to

cover the cases where f and α are complex-valued functions.

**Now let α be any curve [a, b] → R2 . Then α can be expressed in terms of the components
**

(α1 , α2 ) and can be assosciated with the complex valued function

z(t) = α1 (t) + iα2 (t)

**Given any complex-valued function of a complex variable, f say, defined on Im(α) we define
**

the contour integral of f along α, denoted by

intα f (z)dz

by

intα f (z)dz = intba f (z(t))dz(t)

whenever the complex Riemann-Stieltjes integral on the right exists.

(iii) Reversing the direction of the curve changes the sign of the integral.

1584

(iv) The contour integral always exists if α is rectifiable and f is continuous.

(v) If α is piecewise smooth and the countour integral of f along α exists, then

intα f dz = intba f (z(t))z 0 (t)dt

Version: 4 Owner: vypertd Author(s): vypertd

411.7 orientation

**Let α be a rectifiable, Jordan curve in R2 and z0 be a point in R2 − Im(α) and let α have
**

a winding number W [α : z0 ]. Then W [α : z0 ] = ±1; all points inside α will have the same

index and we define the orientation of a Jordan curve α by saying that α is positively

oriented if the index of every point in α is +1 and negatively oriented if it is −1.

Version: 3 Owner: vypertd Author(s): vypertd

**411.8 proof of Weierstrass M-test
**

P

Consider the sequence of partial sums sn = nm=1 fm . Since the sums are finite, each sn is

continuous. Take any p, q ∈ N such that p ≤ q, then, for every x ∈ X, we have

q

X

|sq (x) − sp (x)| = fm (x)

m=p+1

Xq

≤ |fm (x)|

m=p+1

q

X

≤ Mm

m=p+1

P∞

But since n=1 Mn converges, for any > 0 wePcan find an N ∈ N such that, for any

p, q > N and x ∈ X, weP∞ have |sq (x) − sp (x)| ≤ qm=p+1

P∞Mm < . Hence the sequence sn

converges uniformly to n=1 fn , and the function f = n=1 fn is continuous.

Version: 1 Owner: igor Author(s): igor

1585

411.9 unit disk

**The unit disk in the complex plane, denoted ∆, is defined as {z ∈ C : |z| < 1}. The unit
**

circle, denoted ∂∆ or S 1 is the boundary {z ∈ C : |z| = 1} of the unit disk ∆. Every element

z ∈ ∂∆ can be written as z = eiθ for some real value of θ.

Version: 5 Owner: brianbirgen Author(s): brianbirgen

411.10 upper half plane

**The upper half plane in the complex plane, abbreviated UHP, is defined as {z ∈ C : Im(z) >
**

0}.

Version: 4 Owner: brianbirgen Author(s): brianbirgen

411.11 winding number and fundamental group

The winding number is an analytic way to define an explicit isomorphism

W [• : z0 ] : π1 (C \ z0 ) → Z

from the fundamental group of the punctured (at z0 ) complex plane to the group of integers.

Version: 1 Owner: Dr Absentius Author(s): Dr Absentius

1586

Chapter 412

**30B10 – Power series (including
**

lacunary series)

412.1 Euler relation

**Euler’s relation (also known as Euler’s formula) is considered the first bridge between the
**

fields of algebra and geometry, as it relates the exponential function to the trigonometric

sine and cosine functions.

**The goal is to prove
**

eix = cos(x) + i sin(x)

It’s easy to show that

i4n = 1

i4n+1 = i

i4n+2 = −1

i4n+3 = −i

**Now, using the Taylor series expansions of sin x, cos x and ex , we can show that
**

∞ n n

X

ix i x

e =

n=0

n!

X∞

ix x4n ix4n+1 x4n+2 ix4n+3

e = + − −

n=0

(4n)! (4n + 1)! (4n + 2)! (4n + 3)!

Because the series expansion above is absolutely convergent for all x, we can rearrange the

1587

terms of the series as follows

X∞ X∞

ix x2n n x2n+1

e = (−1) + i (−1)n

n=0

(2n)! n=0

(2n + 1)!

eix = cos(x) + i sin(x)

Version: 8 Owner: drini Author(s): drini, fiziko, igor

412.2 analytic

**Let U be a domain in the complex numbers (resp., real numbers). A function f : U −→ C
**

(resp., f : U −→ R) is analytic (resp., real analytic) if f has a Taylor series about each point

x ∈ U that converges to the function f in an open neighborhood of x.

412.2.1 On Analyticity and Holomorphicity

**A complex function is analytic if and only if it is holomorphic. Because of this equivalence,
**

an analytic function in the complex case is often defined to be one that is holomorphic,

instead of one having a Taylor series as above. Although the two definitions are equivalent,

it is not an easy matter to prove their equivalence, and a reader who does not yet have this

result available will have to pay attention as to which definition of analytic is being used.

Version: 4 Owner: djao Author(s): djao

412.3 existence of power series

**In this entry we shall demonstrate the logical equivalence of the holomorphic and analytic
**

concepts. As is the case with so many basic results in complex analysis, the proof of these

facts hinges on the Cauchy integral theorem, and the Cauchy integral formula.

Holomorphic implies analytic.

**Theorem 8. Let U ⊂ C be an open domain that contains the origin, and let f : U → C, be
**

a function such that the complex derivative

f (z + ζ) − f (z)

f 0 (z) = lim

ζ→0 ζ

1588

exists for all z ∈ U. Then, there exists a power series representation

∞

X

f (z) = ak z k , kzk < R, ak ∈ C

k=0

for a sufficiently small radius of convergence R > 0.

**Note: it is just as easy to show the existence of a power series representation around every
**

basepoint in z0 ∈ U; one need only consider the holomorphic function f (z − z0 ).

**Proof. Choose an R > 0 sufficiently small so that the disk kzk 6 R is contained in U. By
**

the Cauchy integral formula we have that

I

1 f (ζ)

f (z) = dζ, kzk < R,

2πi kζk=R ζ − z

**where, as usual, the integration contour is oriented counterclockwise. For every ζ of modulus
**

R, we can expand the integrand as a geometric power series in z, namely

X f (ζ) ∞

f (ζ) f (ζ)/ζ

= = zk , kzk < R.

ζ −z 1 − z/ζ ζ k+1

k=0

**The circle of radius R is a compact set; hence f (ζ) is bounded on it; and hence, the power
**

series above converges uniformly with respect to ζ. Consequently, the order of the infinite

summation and the integration operations can be interchanged. Hence,

∞

X

f (z) = ak z k , kzk < R,

k=0

where I

1 f (ζ)

ak = ,

2πi kζk=R ζ k+1

as desired. QED

**Analytic implies holomorphic.
**

Theorem 9. Let ∞

X

f (z) = an z n , an ∈ C, kzk <

n=0

**be a power series, converging in D = D (0), the open disk of radius > 0 about the origin.
**

Then the complex derivative

f (z + ζ) − f (z)

f 0 (z) = lim

ζ→0 ζ

exists for all z ∈ D, i.e. the function f : D → C is holomorphic.

1589

Note: this theorem generalizes immediately to shifted power series in z − z0 , z0 ∈ C.

**Proof. For every z0 ∈ D, the function f (z) can be recast as a power series centered at z0 .
**

Hence, without loss of generality it suffices to prove the theorem for z = 0. The power series

∞

X

an+1 ζ n , ζ∈D

n=0

**converges, and equals (f (ζ) − f (0))/ζ for ζ 6= 0. Consequently, the complex derivative f 0 (0)
**

exists; indeed it is equal to a1 . QED

Version: 2 Owner: rmilson Author(s): rmilson

**412.4 infinitely-differentiable function that is not ana-
**

lytic

**If f ∈ C∞ , then we can certainly write a Taylor series for f . However, analyticity requires
**

that this Taylor series actually converge (at least across some radius of convergence) to f .

It is not necessary that the power series for f converge to f , as the following example shows.

Let (

e x 6= 0

f (x) = .

0 x=0

Then f ∈ C∞ , and for any n ≥ 0, f (n) (0) = 0 (see below). So the Taylor series for f around

0 is 0; since f (x) > 0 for all x 6= 0, clearly it does not converge to f .

Proof that f (n) (0) = 0

**Let p(x), q(x) ∈ R[x] be polynomials, and define
**

p(x)

g(x) = · f (x).

q(x)

Then, for x 6= 0,

0 (p0 (x) + p(x) x23 )q(x) − q 0 (x)p(x)

g (x) = · e.

q 2 (x)

Computing (e.g. by applying L’Hôpital’s rule), we see that g 0(0) = limx→0 g 0 (x) = 0.

**Define p0 (x) = q0 (x) = 1. Applying the above inductively, we see that we may write
**

f (n) (x) = pqnn(x)

(x)

f (x). So f (n) (0) = 0, as required.

Version: 2 Owner: ariels Author(s): ariels

1590

412.5 power series

**A power series is a series of the form
**

∞

X

ak (x − x0 )k ,

k=0

**with ak , x0 ∈ R or ∈ C. The ak are called the coefficients and x0 the center of the power
**

series. Where it converges it defines a function, which can thus be represented by a power

series. This is what power series are usually used for. Every power series is convergent at

least at x = x0 where it converges to a0 . In addition it is absolutely convergent in the region

{x | |x − x0 | < r}, with

1

r = lim inf p

k→∞ k

|ak |

It is divergent for every x with |x − x0 | > r. For |x − x0 | = r no general predictions can be

made. If r = ∞, the power series converges absolutely for every real or complex x. The real

number r is called the radius of convergence of the power series.

Examples of power series are:

**• Taylor series, for example:
**

∞

X

x xk

e = .

k=0

k!

**• The geometric series:
**

X ∞

1

= xk ,

1−x

k=0

with |x| < 1.

Power series have some important properties:

**• If a power series converges for a z0 ∈ C then it also converges for all z ∈ C with
**

|z − x0 | < |z0 − x0 |.

**• Also, if a power series diverges for some z0 ∈ C then it diverges for all z ∈ C with
**

|z − x0 | > |z0 − x0 |.

**• For |x − x0 | < r Power series can be added by adding coefficients and mulitplied in the
**

obvious way:

∞

X ∞

X

k

ak (x−xo ) · bj (x−x0 )j = a0 b0 +(a0 b1 +a1 b0 )(x−x0 )+(a0 b2 +a1 b1 +a2 b0 )(x−x0 )2 . . . .

k=0 l=0

1591

• (Uniqueness) If two power series are equal and their centers are the same, then their

coefficients must be equal.

**• Power series can be termwise differentiated and integrated. These operations keep the
**

radius of convergence.

Version: 13 Owner: mathwizard Author(s): mathwizard, AxelBoldt

412.6 proof of radius of convergence

**According to Cauchy’s root test a power series is absolutely convergent if
**

p p

lim sup k |ak (x − x0 )k | = |x − x0 | lim sup k |ak | < 1.

k→∞ k→∞

**This is obviously true if
**

1 1

|x − x0 | < lim sup p = p .

k→∞

k

|ak | lim inf k→∞ k |ak |

**In the same way we see that the series is divergent if
**

1

|x − x0 | > p ,

lim inf k→∞ k

|ak |

which means that the right hand side is the radius of convergence of the power series.

**Now from the ratio test we see that the power series is absolutely convergent if
**

ak+1 (x − x0 )k+1 ak+1

lim = |x − x0 | lim < 1.

k→∞ ak (x − x0 )k k→∞ ak

**Again this is true if
**

ak

|x − x0 | < lim .

k→∞ ak+1

**The series is divergent if
**

ak

|x − x0 | > lim ,

k→∞ ak+1

as follows from the ratio test in the same way. So we see that in this way too we can calculate

the radius of convergence.

Version: 1 Owner: mathwizard Author(s): mathwizard

1592

412.7 radius of convergence

**To the power series
**

∞

X

ak (x − x0 )k (412.7.1)

k=0

**there exists a number r ∈ [0, ∞], its radius of convergence, such that the series converges absolutely
**

for all (real or complex) numbers x with |x − x0 | < r and diverges whenever |x − x0 | > r.

(For |x − x0 | = r no general statements can be made, except that there always exists at least

one complex number x with |x − x0 | = r such that the series diverges.)

**The radius of convergence is given by:
**

1

r = lim inf p (412.7.2)

k→∞ k |a |

k

**and can also be computed as
**

ak

r = lim , (412.7.3)

k→∞ ak+1

if this limit exists.

Version: 6 Owner: mathwizard Author(s): mathwizard, AxelBoldt

1593

Chapter 413

**30B50 – Dirichlet series and other
**

series expansions, exponential series

413.1 Dirichlet series

**Let (λn )n≥1 be an increasing sequence of positive real numbers tending to ∞. A Dirichlet
**

series with exponents (λn ) is a series of the form

X

an e−λn z

n

where z and all the an are complex numbers.

**An ordinary Dirichlet series is one having λn = log n for all n. It is written
**

X an

.

nz

The best-known examples are the Riemann zeta fuction (in which an is the constant 1)

and the more general Dirichlet L-series (in which the mapping n 7→ an is multiplicative and

periodic).

When λn = n, the Dirichlet series is just a power series in the variable e−z .

**The following are the basic convergence properties of Dirichlet series. There is nothing
**

profound about their proofs, which can be found in [1] and in various other works on complex

analysis and analytic number theory.

P

Let f (z) = n an e−λn z be a Dirichlet series.

1. If f converges at z = z0 , then f converges uniformly in the region

Re(z − z0 ) ≥ 0 − α ≤ arg(z − z0 ) ≤ α

1594

where α is any real number such that 0 < α < π/2. (Such a region is known as a

“Stoltz angle”.)

**2. Therefore, if f converges at z0 , its sum defines a holomorphic function on the region
**

Re(z) > Re(z0 ), and moreover f (z) → f (z0 ) as z → z0 within any Stoltz angle.

3. f = 0 identically iff all the an are zero.

**So, if f converges somewhere but not everywhere in C, then the domain of its convergence
**

is the region Re(z) > ρ for some real number ρ, which is called the abscissa

P of convergence

of the Dirichlets series. The abscissa of convergence of the series f (z) = n |an |e−λn z , if it

exists, is called the abscissa of absolute convergence of f .

**Now suppose that the coefficients an are all real and ≥ 0. If the series f converges for
**

Re(z) > ρ, and the resulting function admits an analytic extension to a neighbourhood

of ρ, then the series f converges in a neighbourhood of ρ. Consequently, the domain of

convergence of f (unless it is the whole of C) is bounded by a singularity at a point on the

real axis.

**Finally, return to the general case
**

Pofanany complex numbers (an ), but suppose λn = log n, so

f is an ordinary Dirichlet series nz

.

**1. If the sequence (an ) is bounded, then f converges absolutely in the region Re(z) > 1.
**

P

2. If the partial sums ln=k an are bounded, then f converges (not necessarily absolutely)

in the region Re(z) > 0.

Reference:

[1] Serre, J.-P., A Course in Arithmetic, Chapter VI, Springer-Verlag, 1973.

Version: 2 Owner: bbukh Author(s): Larry Hammick

1595

Chapter 414

**30C15 – Zeros of polynomials,
**

rational functions, and other analytic

functions (e.g. zeros of functions with

bounded Dirichlet integral)

414.1 Mason-Stothers theorem

**Mason’s theorem is often described as the polynomial case of the (currently unproven)
**

ABC conjecture.

**Theorem 1 (Mason-Stothers). Let f (z), g(z), h(z) ∈ C[z] be such that f (z) + g(z) = h(z)
**

for all z, and such that f , g, and h are pair-wise relatively prime. Denote the number of

distinct roots of the product f gh(z) by N. Then

max deg{f, g, h} + 1 6 N.

Version: 1 Owner: mathcam Author(s): mathcam

414.2 zeroes of analytic functions are isolated

**The zeroes of a non-constant analytic function on C are isolated. Let f be an analytic
**

function defined in some domain D ⊂ C and let f (z0 ) = 0 for some z0 ∈ D. Because f is

analytic, there is a Taylor series expansion for f around z0 which converges on an open disk

|z − z0 | < R. Write it as f (z) = Σ∞ n

n=k an (z − z0 ) , with ak 6= 0 and k > 0 (ak is the first

non-zero term). One can factor the series so that f (z) = (z − z0 )k Σ∞ n

n=0 an+k (z − z0 ) and

1596

define g(z) = Σ∞ n k

n=0 an+k (z − z0 ) so that f (z) = (z − z0 ) g(z). Observe that g(z) is analytic

on |z − z0 | < R.

**To show that z0 is an isolated zero of f , we must find > 0 so that f is non-zero on
**

0 < |z − z0 | < . It is enough to find > 0 so that g is non-zero on |z − z0 | < by the

relation f (z) = (z − z0 )k g(z). Because g(z) is analytic, it is continuous at z0 . Notice that

g(z0 ) = ak 6= 0, so there exists an > 0 so that for all z with |z − z0 | < it follows that

|g(z) − ak | < |a2k | . This implies that g(z) is non-zero in this set.

Version: 5 Owner: brianbirgen Author(s): brianbirgen

1597

Chapter 415

**30C20 – Conformal mappings of
**

special domains

415.1 automorphisms of unit disk

**All automorphisms of the complex unit disk ∆ = {z ∈ C : |z| < 1} to itself, can be written
**

z−a

in the form fa (z) = eiθ 1−az where a ∈ ∆ and θ ∈ S 1 .

This map sends a to 0, 1/a to ∞ and the unit circle to the unit circle.

Version: 3 Owner: brianbirgen Author(s): brianbirgen

**415.2 unit disk upper half plane conformal equivalence
**

theorem

Theorem: There is a conformal map from ∆, the unit disk, to UHP , the upper half plane.

**Proof: Define f : C → C, f (z) = z−i
**

z+i

1+w

. Notice that f −1 (w) = i 1−w and that f (and therefore

−1

f ) is a Mobius transformation.

−1−i

Notice that f (0) = −1, f (1) = 1−i

1+i

= −i and f (−1) = −1+i = i. By the Mobius circle transformation theore

f takes the real axis to the unit circle. Since f (i) = 0, f maps UHP to ∆ and f −1 : ∆ →

UHP .

Version: 3 Owner: brianbirgen Author(s): brianbirgen

1598

Chapter 416

**30C35 – General theory of conformal
**

mappings

416.1 proof of conformal mapping theorem

**Let D ⊂ C be a domain, and let f : D → C be an analytic function. By identifying the
**

complex plane C with R2 , we can view f as a function from R2 to itself:

f˜(x, y) := (Re f (x + iy), Im f (x + iy)) = (u(x, y), v(x, y))

with u and v real functions. The Jacobian matrix of f˜ is

∂(u, v) ux uy

J(x, y) = = .

∂(x, y) vx vy

As an analytic function, f satisfies the Cauchy-Riemann equations, so that ux = vy and

uy = −vx . At a fixed point z = x + iy ∈ D, we can therefore define a = ux (x, y) = vy (x, y)

and b = uy (x, y) = −vx (x, y). We write (a, b) in polar coordinates as (r cos θ, r sin θ) and get

a b cos θ sin θ

J(x, y) = =r

−b a − sin θ cos θ

**Now we consider two smooth curves through (x, y), which we parametrize by γ1 (t) =
**

(u1 (t), v1 (t)) and γ2 (t) = (u2 (t), v2 (t)). We can choose the parametrization such that

γ1 (0) = γ2 (0) = z. The images of these curves under f˜ are f˜ ◦ γ1 and f˜ ◦ γ2 , respectively,

and their derivatives at t = 0 are

du1

∂(u, v) dγ 1

(f˜ ◦ γ1 ) (0) =

0

(γ1 (0)) · (0) = J(x, y) dvdt

1

∂(x, y) dt dt

**and, similarly, du2
**

(f˜ ◦ γ2 )0 (0) = J(x, y) dt

dv2

dt

1599

by the chain rule. We see that if f 0 (z) 6= 0, f transforms the tangent vectors to γ1 and γ2 at

t = 0 (and therefore in z) by the orthogonal matrix

cos θ sin θ

J/r =

− sin θ cos θ

**and scales them by a factor of r. In particular, the transformation by an orthogonal matrix
**

implies that the angle between the tangent vectors is preserved. Since the determinant of

J/r is 1, the transformation also preserves orientation (the direction of the angle between

the tangent vectors). We conclude that f is a conformal mapping.

Version: 3 Owner: pbruin Author(s): pbruin

1600

Chapter 417

**30C80 – Maximum principle;
**

Schwarz’s lemma, Lindelöf principle,

analogues and generalizations;

subordination

417.1 Schwarz lemma

**Let ∆ = {z : |z| < 1} be the open unit disk in the complex plane C. Let f : ∆ → ∆ be
**

a holomorphic function with f(0)=0. Then |f (z)| ≤ |z| for all z ∈ ∆, and |f 0 (0)| ≤ 1. If

equality |f (z)| = |z| holds for any z 6= 0 or |f 0 (0)| = 1, then f is a rotation: f (z) = az with

|a| = 1.

**This lemma is less celebrated than the bigger guns (such as the Riemann mapping theorem,
**

which it helps prove); however, it is one of the simplest results capturing the “rigidity” of

holomorphic functions. No similar result exists for real functions, of course.

Version: 2 Owner: ariels Author(s): ariels

**417.2 maximum principle
**

Maximum principle Let f : U → R (where U ⊆ Rd ) be a harmonic function. Then f

attains its extremal values on any compact K ⊆ U on the boundary ∂K of K. If f

attains an extremal value anywhere inside int K, then it is constant.

Maximal modulus principle Let f : U → C (where U ⊆ C) be a holomorphic function.

Then |f | attains its maximal value on any compact K ⊆ U on the boundary ∂K of K.

If |f | attains its maximal value anywhere inside int K, then it is constant.

1601

Version: 1 Owner: ariels Author(s): ariels

417.3 proof of Schwarz lemma

**Define g(z) = f (z)/z. Then g : ∆ → C is a holomorphic function. The Schwarz lemma is
**

just an application of the maximal modulus principle to g.

**For any 1 > > 0, by the maximal modulus principle |g| must attain its maximum on the
**

closed disk {z : |z| ≤ 1 − } at its boundary {z : |z| = 1 − }, say at some point z . But then

1

|g(z)| ≤ |g(z )| ≤ 1− for any |z| ≤ 1 − . Taking an infinimum as → 0, we see that values

of g are bounded: |g(z)| ≤ 1.

**Thus |f (z)| ≤ |z|. Additionally, f 0 (0) = g(0), so we see that |f 0 (0)| = |g(0)| ≤ 1. This is the
**

first part of the lemma.

**Now suppose, as per the premise of the second part of the lemma, that |g(w)| = 1 for some
**

w ∈ ∆. For any r > |w|, it must be that |g| attains its maximal modulus (1) inside the

disk {z : |z| ≤ r}, and it follows that g must be constant inside the entire open disk ∆. So

g(z) ⇔ a for a = g(w) of size 1, and f (z) = az, as required.

Version: 2 Owner: ariels Author(s): ariels

1602

Chapter 418

**30D20 – Entire functions, general
**

theory

418.1 Liouville’s theorem

**A bounded entire function is constant. That is, a bounded complex function f : C → C
**

which is holomorphic on the entire complex plane is always a constant function.

**More generally, any holomorphic function f : C → C which satisfies a polynomial bound
**

condition of the form

|f (z)| < c · |z|n

for some c ∈ R, n ∈ Z, and all z ∈ C with |z| sufficiently large is necessarily equal to a

polynomial function.

**Liouville’s theorem is a vivid example of how stringent the holomorphicity condition on
**

a complex function really is. One has only to compare the theorem to the corresponding

statement for real functions (namely, that a bounded differentiable real function is constant,

a patently false statement) to see how much stronger the complex differentiability condition

is compared to real differentiability.

**Applications of Liouville’s theorem include proofs of the fundamental theorem of algebra and
**

of the partial fraction decomposition theorem for rational functions.

Version: 4 Owner: djao Author(s): djao

418.2 Morera’s theorem

Morera’s theorem provides the converse of Cauchy’s integral theorem.

1603

Theorem [1] Suppose G is a region in C, and f : G → C is a continuous function. If for

every closed triangle ∆ in G, we have

int∂∆ f dz = 0,

then f is analytic on G. (Here, ∂∆ is the piecewise linear boundary of ∆.)

**In particular, if for every rectifiable closed curve Γ in G, we have intΓ f dz = 0, then f is
**

analytic on G. Proofs of this can be found in [2, 2].

REFERENCES

1. W. Rudin, Real and complex analysis, 3rd ed., McGraw-Hill Inc., 1987.

2. E. Kreyszig, Advanced Engineering Mathematics, John Wiley & Sons, 1993, 7th ed.

3. R.A. Silverman, Introductory Complex Analysis, Dover Publications, 1972.

Version: 7 Owner: matte Author(s): matte, drini, nerdy2

418.3 entire

A function f : C −→ C is entire if it is holomorphic.

Version: 2 Owner: djao Author(s): djao

418.4 holomorphic

**Let U ⊂ C be a domain in the complex numbers. A function f : U −→ C is holomorphic if
**

f has a complex derivative at every point x in U, i.e. if

f (z) − f (z0 )

lim

z→z0 z − z0

exists for all z0 ∈ U.

Version: 5 Owner: djao Author(s): djao, rmilson

418.5 proof of Liouville’s theorem

**Let f : C → C be a bounded, entire function. Then by Taylor’s Theorem,
**

X∞

1 f (w)

f (z) = cn xn where cn = intΓr n+1 dw

n=0

2πi w

1604

where Γr is the circle of radius r about 0, for r > 0. Then cn can be estimated as

1 f (w) 1 Mr Mr

|cn | 6 length(Γr ) sup n+1 : w ∈ Γr = 2πr n+1 = n

2π w 2π r r

where Mr = sup{|f (w)| : w ∈ Γr }.

**But f is bounded, so there is M such that Mr 6 M for all r. Then |cn | 6 rMn for all n and
**

all r > 0. But since r is arbitrary, this gives cn = 0 whenever n > 0. So f (z) = c0 for all z,

so f is constant.

Version: 2 Owner: Evandar Author(s): Evandar

1605

Chapter 419

**30D30 – Meromorphic functions,
**

general theory

419.1 Casorati-Weierstrass theorem

**Let U ⊂ C be a domain, a ∈ U, and let f : U \ {a} → C be holomorphic. Then a is an
**

essential singularity of f if and only if the image of any punctured neighborhood of a under

f is dense in C.

Version: 2 Owner: pbruin Author(s): pbruin

419.2 Mittag-Leffler’s theorem

**Let G be an open subset of C, let {ak } be a sequence of distinct points in G which has no
**

limit point in G. For each k, let A1k , . . . , Amk k be arbitrary complex coefficients, and define

mk

X Ajk

Sk (z) = .

j=1

(z − ak )j

**Then there exists a meromorphic function f on G whose poles are exactly the points {ak }
**

and such that the singular part of f at ak is Sk (z), for each k.

Version: 1 Owner: Koro Author(s): Koro

1606

419.3 Riemann’s removable singularity theorem

**Let U ⊂ C be a domain, a ∈ U, and let f : U \ {a} be holomorphic. Then a is a
**

removable singularity of f if and only if

lim (z − a)f (z) = 0.

z→a

**In particular, a is a removable singularity of f if f is bounded near a, i.e. if there is a
**

punctured neighborhood V of a and a real number M > 0 such that |f (z)| < M for all

z ∈ V.

Version: 1 Owner: pbruin Author(s): pbruin

419.4 essential singularity

**Let U ⊂ C be a domain, a ∈ U, and let f : U \{a} → C be holomorphic. If the Laurent series
**

expansion of f (z) around a contains infinitely many terms with negative powers of z −a, then

a is said to be an essential singularity of f . Any singularity of f is a removable singularity,

a pole or an essential singularity.

**If a is an essential singularity of f , then the image of any punctured neighborhood of a under
**

f is dense in C (the Casorati-Weierstrass theorem). In fact, an even stronger statement is

true: according to Picard’s theorem, the image of any punctured neighborhood of a is C,

with the possible exception of a single point.

Version: 4 Owner: pbruin Author(s): pbruin

419.5 meromorphic

**Let U ⊂ C be a domain. A function f : U −→ C is meromorphic if f is holomorphic except
**

at an isolated set of poles.

It can be proven that if f is meromorphic then its set of poles does not have an accumulation point.

Version: 2 Owner: djao Author(s): djao

419.6 pole

**Let U ⊂ C be a domain and let a ∈ C. A function f : U −→ C has a pole at a if it can
**

be represented by a Laurent series centered about a with only finitely many negative terms;

1607

that is,

∞

X

f (z) = ck (z − a)k

k=−n

in some nonempty deleted neighborhood of a, for some n ∈ N.

Version: 2 Owner: djao Author(s): djao

419.7 proof of Casorati-Weierstrass theorem

**Assume that a is an essential singularity of f . Let V ⊂ U be a punctured neighborhood of
**

a, and let λ ∈ C. We have to show that λ is a limit point of f (V ). Suppose it is not, then

there is an > 0 such that |f (z) − λ| > for all z ∈ V , and the function

1

g : V → C, z 7→

f (z) − λ

1

is bounded, since |g(z)| = |f (z)−λ| < −1 for all z ∈ V . According to Riemann’s removable singularity theorem

this implies that

S a is a removable singularity of g, so that g can be extended to a holomorphic

function ḡ : V {a} → C. Now

1

f (z) = −λ

ḡ(z)

for z 6= a, and a is either a removable singularity of f (if ḡ(z) 6= 0) or a pole of order n (if ḡ

has a zero of order n at a). This contradicts our assumption that a is an essential singularity,

which means that λ must be a limit point of f (V ). The argument holds for all λ ∈ C, so

f (V ) is dense in C for any punctured neighborhood V of a.

**To prove the converse, assume that f (V ) is dense in C for any punctured neighborhood V
**

of a. If a is a removable singularity, then f is bounded near a, and if a is a pole, f (z) → ∞

as z → a. Either of these possibilities contradicts the assumption that the image of any

punctured neighborhood of a under f is dense in C, so a must be an essential singularity of

f.

Version: 1 Owner: pbruin Author(s): pbruin

**419.8 proof of Riemann’s removable singularity theo-
**

rem

**Suppose that f is holomorphic on U \ {a} and limz→a (z − a)f (z) = 0. Let
**

∞

X

f (z) = ck (z − a)k

k=−∞

1608

be the Laurent series of f centered at a. We will show that ck = 0 for k < 0, so that f can

be holomorphically extended to all of U by defining f (a) = c0 .

**For n ∈ N0 , the residue of (z − a)n f (z) at a is
**

I

n 1

Res((z − a) f (z), a) = lim (z − a)n f (z)dz.

2πi δ→0+ |z−a|=δ

**This is equal to zero, because
**

I

(z − a) f (z)dz ≤ 2πδ max |(z − a)n f (z)|

n

|z−a|=δ

|z−a|=δ

= 2πδ n max |(z − a)f (z)|

|z−a|=δ

**which, by our assumption, goes to zero as δ → 0. Since the residue of (z − a)n f (z) at a is
**

also equal to c−n−1 , the coefficients of all negative powers of z in the Laurent series vanish.

**Conversely, if a is a removable singularity of f , then f can be expanded in a power series
**

centered at a, so that

lim (z − a)f (z) = 0

z→a

because the constant term in the power series of (z − a)f (z) is zero.

A corollary of this theorem is the following: if f is bounded near a, then

|(z − a)f (z)| ≤ |z − a|M

**for some M > 0. This implies that (z − a)f (z) → 0 as z → a, so a is a removable singularity
**

of f .

Version: 1 Owner: pbruin Author(s): pbruin

419.9 residue

**Let U ⊂ C be a domain and let f : U −→ C be a function represented by a Laurent series
**

∞

X

f (z) := ck (z − a)k

k=−∞

**centered about a. The coefficient c−1 of the above Laurent series is called residue of f at a,
**

and denoted Res(f ; a).

Version: 2 Owner: djao Author(s): djao

1609

419.10 simple pole

**A simple pole is a pole of order 1. That is, a meromorphic function f has a simple pole at
**

x0 ∈ C if

a

f (z) = + g(z)

z − x0

where a 6= 0 ∈ C, and g is holomorphic at x0 .

Version: 3 Owner: bwebste Author(s): bwebste

1610

Chapter 420

**30E20 – Integration, integrals of
**

Cauchy type, integral representations

of analytic functions

420.1 Cauchy integral formula

**The formulas. Let D = {z ∈ C : kz − z0 k < R} be an open disk in the complex plane,
**

and let f (z) be a holomorphic 1 function defined on some open domain that contains D and

its boundary. Then, for every z ∈ D we have

I

1 f (ζ)

f (z) = dζ

2πi C ζ − z

I

0 1 f (ζ)

f (z) = dζ

2πi C (ζ − z)2

..

. I

(n) n! f (ζ)

f (z) = dζ

2πi C (ζ − z)n+1

Here C = ∂D is the corresponding circular boundary contour, oriented counterclockwise,

with the most obvious parameterization given by

ζ = z0 + Reit , 0 6 t 6 2π.

**Discussion. The first of the above formulas underscores the “rigidity” of holomoprhic
**

functions. Indeed, the values of the holomorphic function inside a disk D are completely

1

It is necessary to draw a distinction between holomorphic functions (those having a complex derivative)

and analytic functions (those representable by power series). The two concepts are, in fact, equivalent, but

the standard proof of this fact uses the Cauchy Integral Formula with the (apparently) weaker holomorphicity

hypothesis.

1611

specified by its values on the boundary of the disk. The second formula is useful, because it

gives the derivative in terms of an integral, rather than as the outcome of a limit process.

**Generalization. The following technical generalization of the formula is needed for the
**

treatment of removable singularities. Let S be a finite subset of D, and suppose that f (z)

is holomorphic for all z ∈

/ S, but also that f (z) is bounded near all z ∈ S. Then, the above

formulas are valid for all z ∈ D \ S.

**Using the Cauchy residue theorem, one can further generalize the integral formula to the
**

situation where D is any domain and C is any closed rectifiable curve in D; in this case, the

formula becomes I

1 f (ζ)

η(C, z)f (z) = dζ

2πi C ζ − z

where η(C, z) denotes the winding number of C. It is valid for all points z ∈ D \ S which

are not on the curve C.

Version: 19 Owner: djao Author(s): djao, rmilson

420.2 Cauchy integral theorem

**Theorem 10. Let U ⊂ C be an open, simply connected domain, and let f : U → C be a
**

function whose complex derivative, that is

f (w) − f (z)

lim ,

w→z w−z

exists for all z ∈ U. Then, the integral around a every closed contour γ ⊂ U vanishes; in

symbols I

f (z) dz = 0.

γ

**We also have the following, technically important generalization involving removable singularities.
**

Theorem 11. Let U ⊂ C be an open, simply connected domain, and S ⊂ U a finite subset.

Let f : U\S → C be a function whose complex derivative exists for all z ∈ U\S, and that

is bounded near all z ∈ S. Then, the integral around a every closed contour γ ⊂ U\S that

avoids the exceptional points vanishes.

**Cauchy’s theorem is an essential stepping stone in the theory of complex analysis. It is
**

required for the proof of the Cauchy integral formula, which in turn is required for the proof

that the existence of a complex derivative implies a power series representation.

**The original version of the theorem, as stated by Cauchy in the early 1800s, requires that the
**

derivative f 0 (z) exist and be continuous. The existence of f 0 (z) implies the Cauchy-Riemann equations,

1612

which in turn can be restated as the fact that the complex-valued differential f (z) dz is closed.

The original proof makes use of this fact, and calls on Green’s theorem to conclude that the

contour integral vanishes. The proof of Green’s theorem, however, involves an interchange

of order in a double integral, and this can only be justified if the integrand, which involves the

real and imaginary parts of f 0 (z), is assumed to be continuous. To this date, many authors

prove the theorem this way, but erroneously fail to mention the continuity assumption.

In the latter part of the 19th century E. Goursat found a proof of the integral theorem that

merely required that f 0 (z) exist. Continuity of the derivative, as well as the existence of all

higher derivatives, then follows as a consequence of the Cauchy integral formula. Not only

is Goursat’s version a sharper result, but it is also more elementary and self-contained, in

that sense that it is does not require Green’s theorem. Goursat’s argument makes use of

rectangular contour (many authors use triangles though), but the extension to an arbitrary

simply-connected domain is relatively straight-forward.

Theorem 12 (Goursat). Let U be an open domain containing a rectangle

R = {x + iy ∈ C : a 6 x 6 b , c 6 y 6 d}.

If the complex derivative of a function f : U → C exists at all points of U, then the contour

integral of f around the boundary of R vanishes; in symbols

I

f (z) dz = 0.

∂R

Bibliography.

• A. Ahlfors, “Complex Analysis”.

Version: 7 Owner: rmilson Author(s): rmilson

420.3 Cauchy residue theorem

**Let U ⊂ C be a simply connected domain, and suppose f is a complex valued function
**

which is defined and analytic on all but finitely many points a1 , . . . , am of U. Let C be a

closed curve in U which does not intersect any of the ai . Then

m

X

intC f (z) dz = 2πi η(C, ai) Res(f ; ai ),

i=1

where

1 dz

η(C, ai) := intC

2πi z − ai

is the winding number of C about ai , and Res(f ; ai) denotes the residue of f at ai .

Version: 4 Owner: djao Author(s): djao, rmilson

1613

420.4 Gauss’ mean value theorem

**Let Ω be a domain in C and suppose f is an analytic function on Ω. Furthermore, let C be
**

a circle inside Ω with center z0 and radius r. Then f (z0 ) is the mean value of f along C,

that is,

1

f (z0 ) = int2π f (z0 + reiθ )dθ.

2π 0

Version: 7 Owner: Johan Author(s): Johan

420.5 Möbius circle transformation theorem

Möbius transformations always transform circles into circles.

Version: 1 Owner: Johan Author(s): Johan

**420.6 Möbius transformation cross-ratio preservation
**

theorem

A Möbius transformation f : z 7→ w preserves the cross-ratios, i.e.

(w1 − w2 )(w3 − w4 ) (z1 − z2 )(z3 − z4 )

=

(w1 − w4 )(w3 − w2 ) (z1 − z4 )(z3 − z2 )

Version: 3 Owner: Johan Author(s): Johan

420.7 Rouch’s theorem

**Let f, g be analytic on and inside a simple closed curve C. Suppose |f (z)| > |g(z)| on C.
**

Then f and f + g have the same number of zeros inside C.

Version: 2 Owner: Johan Author(s): Johan

1614

420.8 absolute convergence implies convergence for an

infinite product

If an infinite product is absolutely convergent then it is convergent.

Version: 2 Owner: Johan Author(s): Johan

**420.9 absolute convergence of infinite product
**

Q∞ Q∞

An infinite product n=1 (1 + an ) is said to be absolutely convergent if n=1 (1 + |an |) con-

verges.

Version: 4 Owner: mathcam Author(s): mathcam, Johan

420.10 closed curve theorem

Let U ⊂ C be a simply connected domain, and suppose f : U −→ C is holomorphic. Then

intC f (z) dz = 0

for any smooth closed curve C in U.

**More generally, if U is any domain, and C1 and C2 are two homotopic smooth closed curves
**

in U, then

intC1 f (z) dz = intC2 f (z) dz.

for any holomorphic function f : U −→ C.

Version: 3 Owner: djao Author(s): djao

420.11 conformal Möbius circle map theorem

Any conformal map that maps the interior of the unit disc onto itself is a Möbius transformation.

Version: 4 Owner: Johan Author(s): Johan

1615

420.12 conformal mapping

**A mapping f : C 7→ C which preserves the size and orientation of the angles (at z0 ) between
**

any two curves which intersects in a given point z0 is said to be conformal at z0 . A mapping

that is conformal at any point in a domain D is said to be conformal in D.

Version: 4 Owner: Johan Author(s): Johan

420.13 conformal mapping theorem

Let f (z) be analytic in a domain D. Then it is conformal at any point z ∈ D where f 0 (z) 6= 0.

Version: 2 Owner: Johan Author(s): Johan

**420.14 convergence/divergence for an infinite product
**

Q∞

Consider

Qm n=1 pn . We say that this infinite product converges iff the finite products Pm =

p

n=1 n −→ P 6= 0 converge or for at most a finite number of terms pnk = 0 , k = 1, . . . , K.

Otherwise the infinite product is called divergent.

Note: The infinite product vanishes only if a factor is zero.

Version: 6 Owner: Johan Author(s): Johan

420.15 example of conformal mapping

**Consider the four curves A = {t}, B = {t + it}, C = {it} and D = {−t + it}, t ∈ [−10, 10].
**

Suppose there is a mapping f : C 7→ C which maps A to D and B to C. Is f conformal

at z0 = 0? The size of the angles between A and B at the point of intersection z0 = 0 is

preserved, however the orientation is not. Therefore f is not conformal at z0 = 0. Now

suppose there is a function g : C 7→ C which maps A to C and B to D. In this case we

see not only that the size of the angles is preserved, but also the orientation. Therefore g is

conformal at z0 = 0.

Version: 3 Owner: Johan Author(s): Johan

1616

420.16 examples of infinite products

**A classic example is the Riemann zeta function. For Re(z) > 1 we have
**

X∞ Y

1 1

ζ(z) = z

= −z

.

n=1

n p prime

1 − p

**With the help of a Fourier series, or in other ways, one can prove this infinite product
**

expansion of the sine function:

Y∞

z2

sin z = z 1− 2 2 (420.16.1)

n=1

nπ

**where z is an arbitrary complex number. Taking the logarithmic derivative (a frequent move
**

in connection with infinite products) we get a decomposition of the cotangent into partial

fractions: ∞

1 X 1 1

π cot πz = + 2 + . (420.16.2)

z n=1

z+n z−n

The equation (495.2.1), in turn, has some interesting uses, e.g. to get the Taylor expansion

of an Eisenstein series, or to evaluate ζ(2n) for positive integers n.

Version: 1 Owner: mathcam Author(s): Larry Hammick

420.17 link between infinite products and sums

Let ∞

Y

pk

k=1

**be an infinite product such that pk > 0 for all k. Then the infinite product converges if and
**

only if the infinite sum

∞

X

log pk

k=1

converges. Moreover

∞

Y ∞

X

pk = exp log pk .

k=1 k=1

Proof.

**Simply notice that
**

N

Y N

X

pk = exp log pk .

k=1 k=1

1617

If the infinite sum converges then

N

Y N

X ∞

X

lim pk = lim exp log pk = exp log pk

N →∞ N →∞

k=1 k=1 k=1

and also the infinite product converges.

Version: 1 Owner: paolini Author(s): paolini

420.18 proof of Cauchy integral formula

**Let D = {z ∈ C : kz − z0 k < R} be a disk in the complex plane, S ⊂ D a finite subset, and
**

U ⊂ C an open domain that contains the closed disk D̄. Suppose that

• f : U\S → C is holomorphic, and that

• f (z) is bounded on D\S.

Let z ∈ D\S be given, and set

f (ζ) − f (z)

g(ζ) = , ζ ∈ D\S 0,

ζ −z

S

where S 0 = S {z}. Note that g(ζ) is holomorphic and bounded on D\S 0. The second

assertion is true, because

g(ζ) → f 0 (z), as ζ → z.

Therefore, by the Cauchy integral theorem

I

g(ζ) dζ = 0,

C

where C is the counterclockwise circular contour parameterized by

ζ = z0 + Reit , 0 6 t 6 2π.

Hence, I I

f (ζ) f (z)

dζ = dζ. (420.18.1)

C ζ −z C ζ −z

**lemma. If z ∈ C is such that kzk =
**

6 1, then

I (

dζ 0 if kzk > 1

=

kζk=1 ζ −z 2πi if kzk < 1

1618

The proof is fun exercise in elementary integral calculus, an application of the half-angle

trigonometric substitutions.

**Thanks to the Lemma, the right hand side of (495.2.1) evaluates to 2πif (z). Dividing through
**

by 2πi, we obtain I

1 f (ζ)

f (z) = dζ,

2πi C ζ − z

as desired.

Since a circle is a compact set, the defining limit for the derivative

d f (ζ) f (ζ)

=

dz ζ − z (ζ − z)2

**converges uniformly for ζ ∈ ∂D. Thanks to the uniform convergence, the order of the
**

derivative and the integral operations can be interchanged. In this way we obtain the second

formula: I I

0 1 d f (ζ) 1 f (ζ)

f (z) = dζ = dζ.

2πi dz C ζ − z 2πi C (ζ − z)2

Version: 9 Owner: rmilson Author(s): rmilson, stawn

420.19 proof of Cauchy residue theorem

**Being f holomorphic by Cauchy Riemann equations the differential form f (z) dz is closed.
**

So by the lemma about closed differential forms on a simple connected domain we know that

the integral intC f (z) dz is equal to intC 0 f (z) dz if C 0 is any curve which is homotopic to C.

In particular we can consider a curve C 0 which turns around the points aj along small circles

and join these small circles with segments. Since the curve C 0 follows each segment two times

with opposite orientation it is enough to sum the integrals of f around the small circles.

**So letting z = aj + ρeiθ be a parameterization of the curve around the point aj , we have
**

dz = ρieiθ dθ and hence

X

intC f (z) dz = intC 0 f (z) dz = η(C, aj )int∂Bρ (aj ) f (z) dz

j

X

= η(C, aj )int2π iθ iθ

0 f (aj + ρe )ρie dθ

j

**where ρ > 0 is choosen so small that the balls Bρ (aj ) are all disjoint and all contained in the
**

domain U. So by linearity, it is enough to prove that for all j

iint2π iθ iθ

0 f (aj + e )ρe dθ = 2πiRes(f, aj ).

1619

Let now j be fixed and consider now the Laurent series for f in aj :

X

f (z) = ck (z − aj )k

k∈Z

**so that Res(f, aj ) = c−1 . We have
**

X X

int2π

0 f (aj + eiθ

)ρeiθ

dθ = int 2π

0 kc (ρeiθ k

) ρeiθ

dθ = ρk+1

ck int2π

0 e

i(k+1)θ

dθ.

k k

Notice now that if k = −1 we have

ρk+1 ck int2π

0 e

i(k+1)θ

dθ = c−1 int2π

0 dθ = 2πc−1 = 2π Res(f, aj )

**while for k 6= −1 we have
**

2π

ei(k+1)θ

int2π

0 e

i(k+1)θ

dθ = = 0.

i(k + 1) 0

Hence the result follows.

Version: 2 Owner: paolini Author(s): paolini

420.20 proof of Gauss’ mean value theorem

**We can parametrize the circle by letting z = z0 + reiφ . Then dz = ireiφ dφ. Using the
**

Cauchy integral formula we can express f (z0 ) in the following way:

I

1 f (z) 1 f (z0 + reiφ ) iφ 1

f (z0 ) = dz = int2π

0 ire dφ = int2π f (z0 + reiφ )dφ.

2πi C z − z0 2πi reiφ 2π 0

Version: 12 Owner: Johan Author(s): Johan

420.21 proof of Goursat’s theorem

**We argue by contradiction. Set I
**

η= f (z) dz,

∂R

**and suppose that η 6= 0. Divide R into four congruent rectangles R1 , R2 , R3 , R4 (see Figure
**

1), and set I

ηi = f (z) dz.

∂Ri

1620

Figure 1: subdivision of the rectangle contour.

**Now subdivide each of the four sub-rectangles, to get 16 congruent sub-sub-rectangles
**

Ri1 i2 , i1 , i2 = 1 . . . 4, and then continue ad infinitum to obtain a sequence of nested fami-

lies of rectangles Ri1 ...ik , with ηi1 ...ik the values of f (z) integrated along the corresponding

contour.

**Orienting the boundary of R and all the sub-rectangles in the usual counter-clockwise fashion
**

we have

η = η1 + η2 + η3 + η4 ,

and more generally

ηi1 ...ik = ηi1 ...ik 1 + ηi1 ...ik 2 + ηi1 ...ik 3 + ηi1 ...ik 4 .

In as much as the integrals along oppositely oriented line segments cancel, the contributions

from the interior segments cancel, and that is why the right-hand side reduces to the integrals

along the segments at the boundary of the composite rectangle.

**Let j1 ∈ {1, 2, 3, 4} be such that |ηj1 | is the maximum of |ηi |, i = 1, . . . , 4. By the triangle inequality
**

we have

|η1 | + |η2 | + |η3 | + |η4 | > |η|,

and hence

|ηj1 | > 1/4|η|.

Continuing inductively, let jk+1 be such that |ηj1 ...jk jk+1 | is the maximum of |ηj1 ...jk i |, i =

1, . . . , 4. We then have

|ηj1 ...jk jk+1 | > 4−(k+1) |η|. (420.21.1)

**Now the sequence of nested rectangles Rj1 ...jk converges to some point z0 ∈ R; more formally
**

∞

\

{z0 } = Rj1 ...jk .

k=1

**The derivative f 0 (z0 ) is assumed to exist, and hence for every > 0 there exists a k sufficiently
**

large, so that for all z ∈ Rj1 ...jk we have

|f (z) − f 0 (z0 )(z − z0 )| 6 |z − z0 |.

Now we make use of the following.

Lemma 9. Let Q ⊂ C be a rectangle, let a, b ∈ C, and let f (z) be a continuous, complex

valued function defined and bounded in a domain containing Q. Then,

I

(az + b)dz = 0

∂Q

I

f (z) 6 MP,

∂Q

where M is an upper bound for |f (z)| and where P is the length of ∂Q.

1621

The first of these assertions follows by the fundamental theorem of calculus; after all the

function az + b has an anti-derivative. The second assertion follows from the fact that the

absolute value of an integral is smaller than the integral of the absolute value of the integrand

— a standard result in integration theory.

Using the lemma and the fact that the perimeter of a rectangle is greater than its diameter

we infer that for every > 0 there exists a k sufficiently large that

I

ηj1 ...jk = f (z) dz 6 |∂Rj1 ...jk |2 = 4−k |∂R|2 .

∂Rj ...j

1 k

**where |∂R| denotes the length of perimeter of the rectangle R. This contradicts the earlier
**

estimate (419.21.1). Therefore η = 0.

Version: 10 Owner: rmilson Author(s): rmilson

420.22 proof of Möbius circle transformation theorem

Case 1: f (z) = az + b.

**Case 1a: The points on |z − C| = R can be written as z = C + Reiθ . They are mapped to
**

the points w = aC + b + aReiθ which all lie on the circle |w − (aC + b)| = |a|R.

iθ

Case 1b: The line Re(eiθ z) = k are mapped to the line Re e aw = k + Re ab .

Case 2: f (z) = z1 .

**Case 2a: Consider a circle passing through the origin. This can be written as |z − C| = |C|.
**

This circle is mapped to the line Re(Cw) = 21 which does not pass through the origin. To

show this, write z = C + |C|eiθ . w = 1z = C+|C|e

1

iθ .

1 1 C C

Re(Cw) = (Cw + Cw) = iθ

+

2 2 C + |C|e C + |C|e−iθ

1 C C eiθ C/|C| 1 C |C|eiθ 1

= + = + =

2 C + |C|eiθ C + |C|e−iθ eiθ C/|C| 2 C + |C|eiθ |C|eiθ + C 2

Case 2b: Consider the line which does not pass through the origin. This can be written as

Re(az) = 1 for a 6= 0. Then az + az = 2 which is mapped to wa + wa = 2. This is simplified

as aw + aw = 2ww which becomes (w − a/2)(w − a/2) = aa/4 or w − a2 = |a| 2