Educative Commentary on

JEE 2003 MATHEMATICS PAPERS

Note: This material is essentially a supplement to the book ‘Educative JEE

Mathematics’ by K. D. Joshi, published by Universities Press, Hyderabad.
All the references to chapters and exercises are to those in this book. The
text of the questions is taken from the website http://www.iitb.ac.in. The
only change made is that in the Screening Paper the questions are numbered
serially from 1 to 28 instead of the original numbering from 57 to 84. (In the
complete Screeing Paper, the Physics questions were numbered from 1 to 28,
the Chemistry questions from 29 to 56 and the Mathematics questions from
57 to 84.)

In the years prior to 2003, only the JEE question papers were available
to the general public on payment of a nominal charge. In 2003, a welcome
change has taken place. Not only the questions in the Screening Papers but
their answers too are available on the internet. Even the solutions of the
Main Paper are availble. This obviates the need to duplicate them here.
Instead, we give the questions and some comments as to how they can be
tackled (including, in some cases, how they should not be tackled!). This will
illustrate most of the tips given in Chapter 24.
Throughout this book, multiple choice type questions have generally been
converted to the conventional form where you have to find the answer rather
than select it from the given choices. Here, however, we give the questions
exactly as they appeared in the JEE. This is done to illustrate that in multiple
choice questions, the form of the answers provided sometimes gives a clue to
the solution of the problem.
For the questions in the Screening Paper, we first give the answer. This is
followed by how the answer is arrived at and then by comments pointing out
if there is something remarkable about the problem and alternate solutions, if
any, or about its resemblance to some other problem and finally, its possible
extensions and generalisations.
For the problems in the Main Paper, on the other hand, we first give an
analysis of the problem, beginning with possible lines of attack. One of these
lines is then followed till we reach the solution. The discussion given is of
an educative type intended to help you cope with other problems. When
you write the solution in the examination, it is hardly necessary to reveal
how you analysed the problem. It is quite all right to write the major steps
with brief justifications as elaborated in Comment No. 9 of Chapter 24. The
solutions available on the internet can be taken as a good guide to get an
idea of which steps are expected from the candidates.

2
 SCREENING PAPER OF JEE 2003

Q. 1 The centre of the circle inscribed in the square determined by the two
pairs of lines x2 − 8x + 12 = 0 and y 2 − 14y + 45 = 0 is

(A) (7, 3) (B) (7, 4) (C) (4, 7) (D) (3, 7)

Answer and Comments: (C). Factoring the expressions, the sides
lie along the lines x = 2, x = 6, y = 5 and y = 9. The key idea is that
the centre of the inscribed circle in a square is the midpoint of each
diagonal. So it suffices to find two diagonally opposite vertices. One
such pair is (2, 5) and (6, 9). Hence the centre is at (4, 7).
If the pairs of lines were not parallel to the axes, then one would have
to factorise each equation analogously to that in the Main Problem of
Chapter 9. Once you factor one equation, factoring the other becomes
easier because the slopes of the lines represented by its factors are
already known. In the present problem, all this task is eliminated
because the lines are parallel to the axes. It is also easy to find the
perpendicular distance between parallel sides. It is simply the difference
between the roots of each quadratic. Both these differences ought to be
equal as otherwise the rectangle would not be a square. In the present
problem, these differences are 6−2 and 9−5, both of which are equal to
4. Although this is not relevant to the present problem, it is a good idea
to quickly check such things, as they often alert you in case you have
made a numerical slip. Also, if the problem had asked not the centre
but the area of the inscribed circle, then the perpendicular distance
between parallel sides would be directly relevant, because that would
be precisely the diameter.

Q. 2 The coefficient of t24 in the expansion of (1 + t2 )12 (1 + t12 )(1 + t24 ) is

12 12 12 12
(A) C6 + 3 (B) C6 (C) C6 + 1 (D) C6 + 2

3
 Answer and Comments: (D). Evidently, it is foolish to expand the
entire expression. The product of the last two factors is (1 + t12 + t24 +
t36 ). When this is multiplied by some polynomial in t, the only way
t24 will come in the product is if the terms 1, t12 and t24 are multiplied
by t24 , t12 and 1 respectively. So we must find the coefficients of these
powers of t in the first factor, viz., (1 + t2 )12 . By the binomial theorem,
these coefficients are 1, 12C6 and 1 respectively. Hence the desired
coefficient is 12C6 + 2.
The key idea of the problem, viz. the binomial theorem is easy to
get. The only catch is that one should avoid the temptation to expand
the first factor before multiplying the other two. As the expression
deals with only even powers of t, the substitution u = t2 will reduce
the problem to that of finding the coefficient of u6 in the product (1 +
u)12 (1 + u6 )(1 + u12 ). But this does not simplify the work materially.

Q. 3 If the angles of a triangle are in the ratio 4 : 1 : 1, then the longest side
and the perimeter are in the ratio
√ √ √
(A) 3:2+ 3 (B) 1 : 6 (C) 1 : 2 + 3 (D) 2 : 3
Answer and Comments: (A). Clearly, this is a problem on the solu-
tion of triangles (Chapter 11). The data does not determine the triangle
uniquely. It determines it only upto similairity. But that is sufficient
information to find the desired ratio since that ratio is also invariant
under similarity. The key idea of the problem is that the sides of a
triangle are proportional to the sines of their opposite angles.
Let the triangle be ABC. Since A + B + C = 180◦ , from the
data A, B, C must be 120◦ , 30◦ and 30 √
◦
respectively.
√
The desired ratio
a sin 120◦ 3/2 3
is a+b+c = sin 120◦ +sin 30◦ +sin 30◦ = √3/2+1/2+1/2 = 2+√3 . (A student who
finds the angles correctly from the data but does not immediately √ know
what to do with them can take a clue from the presence of 3 in some
of the given alternative answers. Evidently, the sines or the cosines of
the angles are involved. Of course, a student who is perceptive enough
to take this hint will probably not need it!)
The data is a twisted way of specifying the angles of the triangle.
Once the angles are known, the triangle is determined upto similarity.
So anything which depends only on the relative proportions of the

4
 sides (and not on their actual lengths) can be determined from the
angles. For example, we cannot determine either the circumradius R
or the inradius r. But we can easily determine their ratio R/r from
Equation (4) of Chapter 11. For the present triangle it comes out to
be 4 sin 60◦1sin2 15◦ = √3(1−cos
1
30◦ )
= √3(1−1√3/2) = 2√23−3 .
a
However, we cannot go the other way, i.e. from the ratio a+b+c
alone we cannot determine the angles A, B, C. We need one more
piece of information. For example, if we are given that b = c, then
we have B = C and hence A + 2B = 180◦ which implies√sin A =
sin A
sin 2B = 2 sin B cos B. Combining this with sin A+2 sin B
= 2+√3 3 gives
√ √
cos B
cos B+1
= 2+√3 3 and hence cos B = 23 . This determines B (and C) as
a
30 each and hence A = 120◦ . Similarly, the two ratios a+b+c
◦
and Rr
would together determine A, B, C. But the calculations involved are
rather laborious.
R1
Q. 4 For natural numbers m and n, if I(m, n) = tm (1+t)n dt, then I(m, n)
0
equals

n m
(A) I(m + 1, n − 1) (B) I(m + 1, n − 1)
m+ n
1 n+1
2 n
(C) − I(m + 1, n − 1)
m+ n
1 m+1
2 n
(D) +m I(m + 1, n − 1)
m+1 m+1

Answer and Comments: (C). This is very analogous to the Main

Problem of Chapter 18. Integration by parts gives
1
tm+1 n n R 1 m+1
I(m, n) = m+1
(1 + t)  − m+1 0 t (1 + t)n−1 dt
0
2n n
= m+1
− m+1 I(m + 1, n − 1)

In this problem, familiarity with a similar problem enables you to

see the key idea, viz., integration by parts, quickly. Even otherwise, a
perceptive student can think of it because in the original question, all
the four possible answers involved factors with m + 1 or n + 1 in the
denominators.

5
Q. 5 The system of equations

x + ay = 0
y + az = 0
z + ax = 0

has infinitely many solutions for

(A) a = 1 (B) a = 0 (C) a = −1 (D) no real value of a

Answer and Comments: (C). Clearly, this is a problem about the
number of solutions to a homogeneous system of three linear equations
in three unknowns x, y, z. Because of homogeneity, once the system
has one non-trivial solution, say (x0 , y0 , z0 ), then all its multiples, viz.,
(λx0 , λy0 , λz0 ) are also solutions for every real λ. Hence the problem
is equivalent to asking the values of a for which the given system has
one non-trivial solution. By Theorem 8 of Chapter 3, this will be the
case if and only if the determinant  of the coefficients of the system,
 1 a 0 
 
viz. the determinant  0 1 a  vanishes. By a direct computation,
 
 a 0 1 
 

the determinant is a3 + 1 which has only one real root, viz. a = −1.
So this is the only real value for which the system has infinitely many
solutions.
If you don’t know or are not sure of the determinant criterion for
existence of solutions (which is the key idea of the problem), an ad-hoc
solution is possible. Eliminating y and z in the given three equations
gives (a3 + 1)x = 0. Because of the cyclic symmetry of the system, we
also get (a3 + 1)y = 0 and (a3 + 1)z = 0. Now in a non-trivial solution,
at least one of x, y and z is non-zero. So a3 + 1 must vanish, giving
a = −1.

Q. 6 The slopes of the focal chords of the parabola y 2 = 16x which are tan-
gents to the circle (x − 6)2 + y 2 = 2 are
1 1
(A) 1, −1 (B) , −2 (C) − , 2 (D) 2, −2
2 2
Answer and Comments: (A). The focus is at (4, 0) and the equa-
tion of a focal chord with slope m is therefore, y = m(x − 4). The

6
 problem is to find those values of m for which this line touches the
given circle. This can be done using the condition for tangency. But
instead of taxing your memory unnecessarily, you might as well equate
the perpendicular distance of the centre (which is at (6, 0))
 with the
√  −2m  √
radius (which is 2). This gives the equation  √  = 2, i.e.
 
 1 + m2 
4m2 = 2(1 + m2 ) which gives m2 = 1. Hence the slopes are 1 and −1.
The very fact that the equation of the circle is given in a form
from which the centre and the radius can be read off instantaneously,
suggests that it is better not to apply the condition for tangency, and
instead, directly equate the radius with the perpendicular distance of
the centre from the focal chord. Note also that since the parabola as
well as the circle are symmetric about the x-axis, so will be the two
desired focal chords. Hence their slopes must be negatives of each other.
Awareness of this fact can alert you against a possible numerical slip.
It can also help you quickly eliminate some of the false alternatives.
Q. 7 Let A, B, and C be three events and Ā, B̄ and C̄ be their corresponding
complementary events. If the probabilities of the events B, A ∩ B ∩ C̄
3 1 1
and Ā ∩ B ∩ C̄ are , and respectively, then the eprobability of the
4 3 3
event B ∩ C is
1 1 1 1
(A) (B) (C) (D)
9 12 15 18
Answer and Comments: (B). The event B is the disjunction of the
mutually exclusive events B ∩ C and B ∩ C̄. Since P (B) is given, we
can find P (B ∩ C) by subtracting P (B ∩ C̄) from it. But B ∩ C̄ is the
disjunction of the mutually exclusive events A ∩ B ∩ C̄ and Ā ∩ B ∩ C̄.
So P (B ∩ C̄) = P (A ∩ B ∩ C̄) + P (Ā ∩ B ∩ C̄) = 13 + 31 = 23 . Hence
P (B ∩ C) = P (B) − P (B ∩ C̄) = 34 − 32 = 12
1
.
The key idea is that every event E is the disjunction of the two
mutually exclusive events E ∩ F and E ∩ F̄ where F is any event. It
is easy to see this through Venn diagrams. But it is equally clear by
common sense. It is hardly necessary to tax your memory by remem-
bering formulas like (5) in Chapter 22 (whose special case is being used
in this problem twice).
Q. 8 The area of the region in the first quadrant that is bounded by the

7
 √
curves y = x, x = 2y + 3 and the x-axis is
√ 34
(A) 2 3 (B) 18 (C) 9 (D)
3
Answer and Comments: (C). First make a rough sketch of the
region.

y
2
x=y

(9, 3)

O (3, 0) x
x=2y+3

The key idea is to recognise that a horizontal slicing is preferable

since a vertical slicing will require that we decompose the region into
√
two parts by the vertical line x = 3. Rewriting the equation y = x
as x = y 2, the region is bounded on the left by x = y 2 and on the right
by x = 2y + 3. Also 0 ≤ y ≤ 3 for the region. Hence its area is the
R3 y 3 3

integral (2y + 3 − y 2) dy = (y 2 + 3y − )
3 0
 = 9 + 9 − 9 = 9.
0

x−4 y−2 z−k

Q. 9 The line = = lies completely in the plane 2x − 4y +
1 1 2
z = 7 for

(A) no value of k (B) k = 7 (C) k = 1 (D) k = −7

Answer and Comments: (B). The key idea is to notice that the way
the line is given it is very easy to express a typical point on it in terms
of a parameter, say r. Equating each ratio with r we have x = r + 4,
y = r + 2 and z = 2r + k for a point (x, y, z) on the line. This point will
lie on the given plane if and only if 2(r + 4) − 4(r + 2) + (2r + k) = 7.
This gives k = 7.
Here the coefficient of r in the equation we got was 0. If it were non-
zero, then there is no value of k for which the line will lie completely
in the given plane. In that case, for every k the equation will give the
point of intersection of the given line with the given plane. In fact,

8
this shows that one cannot get the answer simply by checking whether
the point (4, 2, k) (which surely lies on the given line) lies in the given
plane. Doing so does give k = 7. But it leaves open the possibility that
the line may intersect the plane at (4, 2, 7). However, if all the multiple
choices given are of the form k = ..., then this approach would work.
(The papersetters were probably aware of this and to prevent such an
unwarranted short cut have given ‘no value of k’ as a possible choice.
The disturbing possibility remains, however, that a student who misses
this point can still mark the correct answer by merely checking for
which k the point (4, 2, k) lies on the plane.)
However, after noting that the point (4, 2, 7) is common to the
given line and the given plane, there are other ways to determine if the
line lies entirely in the plane. One method is to observe that a line lies
completely in a plane as soon as any two points on it do so. So, if we
can find some point on the line, other than (4, 2, 7), which lies in the
given plane, we are through. One such point is obtained by taking each
ratio, i.e. r to be 1. This gives (5, 3, 9) as another point on the line and
it is easy to check directly that its coordinates satisfy the equation of
the plane.
Yet another method is to write the given line as the intersection
of any two planes, e.g. the planes x − 4 = y − 2 and the plane 2(y −
2) = z − 7. These two equations along with the equation of the given
plane give a non-homogeneous system of three linear equations in three
unknowns x, y, z. We already know that it has one solution. To say
that the line lies completely in the plane is equivalent to saying that the
corresponding homogeneous system has infinitely many solutions. And
this can be done by the determinant criterion applied in Q. 5 above. In
fact, this method is better when the line is given as the intersection of
two planes and not in the parametric form. Of course, in such a case,
if one wants, one can first write the line in a parametric form as was
done in the solution to Exercise (21.16).
One more method is provided by vectors. From the data, it is
clear that the vector, say u = i + j + 2k, is parallel to the given line.
Also, the vector, say v = 2i − 4j + k (obtained from the equation of the
plane), is normal to the given plane. Taking dot product, it is easily
checked that u and v are perpendicular to each other. Hence the line
is perpendicular to the normal to the plane. Therefore it is either in

9
 the plane or parallel to it. The second possibility is already excluded.
" # " #
α 0 1 0
Q. 10 If A = and B = are two matrices, then A2 = B is
1 1 5 1
true for

(A) α = −1 (B) α = 1 (C) α = 4 (D) no real value of α

Answer " and Comments:
# (D). By a direct compuation, A2 comes
α2 0
out as . If this equals B then the respective entries must
α+1 1
be equal. In particular, α2 = 1 and α + 1 = 5. But these two equations
have no common solution.
As only 2 × 2 matrices were involved in this problem, it was not
time consuming to compute A2 explicitly. For matrices of higher order,
this could be laborious. Sometimes, short-cuts are available to show
that one matrix cannot be the square of another. Suppose for example,
that A2 = B. Then taking determinants of both the sides and applying
the product rule, we have |A|2 = |B|. So, if |A|2 6= |B|, then we know
A2 6= B even without computing A2 . In fact, if |B| is negative, then
no real matrix A can be a ‘square root’ of B. However, this is only a
one-way conclusion. Even if |A|2 = |B| it does not follow that A2 = B.
In the present problem, for example, |A| = α and |B| = 1. So, if at all
A2 = B, then we must have α2 = 1 and hence α = ±1. But for neither
of these values does A2 equal B.

Q. 11 The minimum value of f (x) = x2 +2bx+2c2 is more than the maximum

value of g(x) = −x2 − 2cx + b2 , x being real, for

√ √
(A) |c|
√ > |b| 2 (B) 0 < c < b 2
(C) b 2 < c < 0 (D) no values of b and c

Answer and Comments: (A). The most straightforward way is to

first determine the minimum of f (x) and the maximum of g(x). This
can be done by calculus. But since the functions are quadratic poly-
nomials, we may as well complete the squares as in Comment No. 11
of Chapter 3. Thus, f (x) = (x + b)2 + 2c2 − b2 whose minimum is
2c2 − b2 (occurring when x = −b, although that is not needed here).

10
 Similarly, since g(x) = −(x + c)2 + b2 + c2 , it has b2 + c2 as its maxi-
mum. The question now reduces to deciding the conditions under which
2c2 − b2 > b2 + c2 , or equivalently, c2 > 2b2 . Taking square roots gives
the answer. (Care has to be taken to write the square root of b2 as |b|
and not either as b or as −b because nothing has been given about the
sign of b. Similarly for c.)

Q. 12 If y(t) is a solution of the differential equation

dy
(1 + t) − ty = 1; y(0) = −1
dt
then the value of y(t) at t = 1 is
1 1 1 1
(A) − (B) e − (C) (D) e +
2 2 2 2
Answer and Comments: (A). This is one of those problems which
are only superficially in the multiple choice format. That is, there is
no way to eliminate even some of the wrong answers. The only way
to arrive at the correct answer is by actually solving the problem as a
conventional one. The given differential equation is linear and of first
order. To solve it we first cast into the standard form
dy −t 1
+ y=
dt 1 + t 1+t
(cf. Equation R(31) in Chapter 19). Then an integrating factor is eF (t)
−t R −1−t+1 R 1
where F (t) = 1+t dt = 1+t
dt = (−1 + 1+t ) dt = −t + ln(1 + t).
F (t) −t+ln(1+t) −t
Hence e =e = e (1 + t). Multiplying the equation by the
integrating factor gives
dy
e−t (1 + t) − te−t y = e−t
dt
(Of course, we could have also got this directly from the original
equation if we could notice that e−t is an integrating factor for it. But
not everybody is sharp enough to do so.)
d
Rewriting the new equations as [(1 + t)e−t y] = e−t and integrating
dt
both the sides w.r.t. t gives the general solution as (1+t)e−t y = −e−t +
C where C is an arbitrary constant. The initial condition y(0) = −1

11
 determines C as 0. (This is a little unfortunate because a candidate
who simply misses the arbitrary constant can also luckily get the correct
answer.) Putting t = 1 now gives y = − 21 .
Note that although we have used the formula given in Comment No.
12 of Chapter 19 for the integrating factor of a linear differential equa-
tion, we have not used the formula (38) there for the general solution
of such a d.e. Once we get the integrating factor, the rest of the work
is very easy. It is not a good idea to tax your memory by remembering
formulas which can be easily derived fresh.

Q. 13 Let f (x) be a differentiable function with f ′ (1) = 4 and f ′ (2) = 6,

where f ′ (c) is the derivative of f (x) at x = c. Then the limit of
f (2 + 2h + h2 ) − f (2)
, as h → 0,
f (1 + h − h2 ) − f (1)
3
(A) may not exist (B) equals 3 (C) equals (D) equals −3
2
f (2 + h) − f (2)
Answer and Comments: (B). The limits of the ratios
h
f (1 + h) − f (1)
and are, by definition, the derivatives f (2) and f ′ (1)
′
h
respectively. So it is tempting to do the problem by dividing both the
numerator and the denominator by h. The trouble is that the numer-
ator is not f (2 + h) − f (2) but rather f (2 + 2h + h2 ) − f (2). Similarly,
the denominator is f (1 + h − h2 ) − f (1) and not f (1 + h) − f (1). The
key idea is to note that as h → 0, h2 is insignificant in comparison to h
and so we may as well drop it. (Similar considerations were applied for
finding a limit at the end of Comment No. 14 of Chapter 15.) Then

f (2 + 2h + h2 ) − f (2) f (2 + 2h) − f (2)

lim = lim
h→0 h h→0 h
f (2 + 2h) − f (2)
= 2 lim
h→0 2h
= 2f ′ (2) = 12

f (2 + 2h) − f (2)
Note the slight hitch in the last but one step. lim
h→0 h
is not quite the derivative of f at 2. It would be so had there been 2h
instead of h in the denominator for then one can replace 2h by u which

12
also tends to 0 as h tends to 0. So we multiply both the numerator and
the denominator by 2 and then take the limit.
If one wants, this argument can be replaced by a more rigorous one.
f (2 + 2h + h2 ) − f (2) f (2 + 2h + h2 ) − f (2) 2h + h2
We rewrite as ×
h 2h + h2 h
which is valid because because if h lies in a sufficiently small deleted
neighbourhood of 0, then h2 + 2h 6= 0. Now as h → 0, h2 + 2h also
tends to 0 and so the first ratio tends to f ′ (2). The second ratio tends
to 2 as h → 0.
f (1 + h − h2 ) − f (1) f (1 + h) − f (1)
Similarly, lim = lim = f ′ (1) = 4.
h→0 h h→0 h
Hence the answer to the problem is 12 4
= 3.
More generally, the argument shows that if A and B are non-zero
constants and θ(h), ψ(h) are functions of h which tend to 0 more rapidly
f (a + Ah + θ(h)) ′ (a)
than h as h → 0, then → Af
Bf ′ (b)
as h → 0. We advise
f (b + Bh + ψ(h))
you against remembering this as a readymade result. Just remember
the reasoning behind it. If A = 0 and B 6= 0 then this limit is 0. But
if A 6= 0 and B = 0, then this limit may not exist. In case both A and
B are 0, the fate of the limit will depend upon the functions θ(h) and
ψ(h), and in particular on their relative rate of tending to 0 as h tends
to 0.
Note that in the given problem, the expression is of the form 00
as h → 0. Hence it is tempting to get the answer using l’Hôpital’s
rule. To do so, we must first regard both the numerator and the de-
nominator as functions of h and take their derivatives w.r.t. h. Us-
ing the chain rule, for the numerator this derivative is f ′ (2 + 2h +
h2 )(2 + 2h) while for the denominator it is f ′ (1 + h − h2 )(1 − 2h).
f ′ (2 + 2h + h2 )(2 + 2h)
So, in case the ratio ′ tends to some limit L,
f (1 + h − h2 )(1 − 2h)
then by l’Hôpital’s rule, L will also be the limit of the original ratio.
2 + 2h
Now, evidently lim = 2. But we cannot so readily conclude that
h→0 1 − 2h
f ′ (2 + 2h + h2 ) f ′ (2)
lim ′ = . For this we need the continuity of f ′ (x)
h→0 f (1 + h − h2 ) f ′ (1)
at 1 and at 2. As this is not given to be the case, the solution based on
l’Hôpital’s rule is not quite adequate. But since the question only asks
for the answer and not the method, there is no harm in applying it.

13
 (Yet another instance where in a multiple choice question, a candidate
can score even though his method is not correct.)

Q. 14 If n is a non-zero integer and

sin(nx)((a − n)nx − tan(x))
lim =0
x→0 x2
then, a equals
1 n+1 1
(A) (B) 0 (C) (D) n +
n n n
Answer and Comments: (D). The key idea is that for x near 0, sin nx
is nearly nx while tan x is nearly x. So, if we rewrite the expression as
sin(nx) (a − n)nx − tan(x)
× then the first factor tends to the limit n
x x
as x → 0. As n is given to be non-zero, the limit of the second factor
must be 0. But either by replacing tan(x) by x or by using l’Hôpital’s
rule, the limit of the second factor is (a − n)n − 1. Setting this equal
to 0 gives a = n + n1 .

x2 + x + 2
Q. 15 The range of the function f (x) = , x ∈ (−∞, ∞) is
x2 + x + 1
(A) [1, ∞) (B) (1, 11/7) (C) (1, 7/3] (D) [1, 7/5]
1
Answer and Comments: (C). The simplification f (x) = 1+
1 + x + x2
suggests itself and reduces the problem to that of findind the range of
1
the function 2 = g(x), say. For this, a variety of methods is
x +x+1
available as shown in Chapters 6 and 10. In the present case, since the
numerator of g(x) is a constant, the problem can be done by study-
ing the behaviour of its denominator, viz. x2 + x + 1 as x varies over
(−∞, ∞). Rewriting x2 + x + 1 as (x + 12 )2 + 43 we see that it is never
zero. In fact, it is always positive and its minimum value is 34 (occur-
ring when x = − 12 ). There is no maximum for the denominator. As x
increases from − 12 to ∞, x2 +x+1 increases from 43 to ∞. So, the range
of the denominator of g(x) is [3/4, ∞). Correspondingly, the range of
g(x) is (0, 34 ]. Since f (x) = 1 + g(x), the range of the original function
f (x) is obtained by shifting this to the right by 1. This gives, (1, 73 ] as
the answer.

14
 Although not particularly recommended, the problem can also be
done by other methods. In each method, transforming the problem
from f (x) to g(x) simplifies the matters considerably. We thus have an
excellent illustration of the remarks in Comment No. 5(iv) of Chapter
24. The direct method of finding the range of g(x) is to identify all
1
those real numbers p for which the equation 2 = p has real
x +x+1
roots. Clearly for p = 0, there is no solution. For p 6= 0, the problem is
equivalent to determining whether the quadratic px2 + px + (p − 1) = 0
has a real solution. The discriminant criterion gives p2 ≥ 4p(p − 1).
Depending on the sign of p, this reduces to two cases : (i) p > 0 and
p ≥ 4(p − 1) and (ii) p < 0 and p ≤ 4(p − 1). The first possibility gives
0 < p ≤ 43 , while the second possibility can hold for no value of p. Thus
the range of g(x) is (0, 43 ] and hence that of f (x) is (1, 37 ].
Finally, there is the calculus method in which we consider g ′ (x) =
2x + 1
− 2 for all real x. The denominator is always positive.
(x + x + 1)2
Hence g ′ vanishes only at x = − 12 , is positive for x < − 12 and negative
for x > − 12 . Hence the function g(x) is strictly increasing on (−∞, − 21 )
and strictly decreasing on (− 12 , ∞). Therefore it has a global maximum
at x = − 12 . By direct computation, g(− 21 ) = 43 . Also clearly, g(x) → 0
as x → ±∞. But g(x) never assumes the value 0. Putting it all
together, the range of g(x) is the interval (0, 34 ]. Shifting it by 1 to the
right gives the range of f as (1, 73 ].

Q. 16 The volume of the parallelopiped formed by vectors î + aĵ + k̂, ĵ + ak̂

and aî + k̂, a > 0 is minimum when the value of a is
√ 1 1
(A) 3 (B) 3 (C) √ (D)
3 3
Answer and Comments: (C). Another straightforward problem,
where the best course of action is to first compute the volume, say
V , of the paralleopiped and then  to minimise
 this as a function of a.
 1 a 1 
 
A direct calculation gives V =  0 1 a  = a3 − a + 1 = f (a) (say).
 
 a 0 1 
 

We have to minimise f (a) for a ∈ (0, ∞). Since f ′ (a) = 3a2 − 1, the
only (positive) critical point is a = √13 . Since the question is a multiple

15
 choice one, one may as well stop here. But, for an honest solution, we
verify that f ′ (a) < 0 for 0 < a < √13 while f ′ (a) > 0 for a > √13 . Hence
f attains its minimum on (0, ∞) at a = √13 .
x
Q. 17 The function f [0, ∞) → [0, ∞), defined by f (x) = , is
1+x

(A) one-to-one and onto

(B) one-to-one but NOT onto
(C) onto but NOT one-to-one
(D) NEITHER one-to-one NOR onto

Answer and Comments: (B). Clearly, f (x) can never be 1 and hence
f is not onto. This eliminates the choices (A) and (C). To check whether
f is one-to-one, one can proceed directly from definition. Thus suppose
a b
f (a) = f (b) for some a, b ∈ [0, ∞). Then 1+a = 1+b which implies a = b
directly. Hence f is one-to-one.
The relatively simple nature of the function f makes it eassy to
x y
obtain the inverse function. Solving y = 1+x for x gives x = 1−y .
−1 −1 x
Hence f is the function defined by f (x) = 1−x . However, this is
not defined at x = 1. So the original function can’t be onto, even
though the existence of the inverse implies that it is one-to-one.

Q. 18 If z is a complex number satisfying |z| = 1 and z 6= −1, then the real

z−1
part of w = is
z+1
√
1 −1 2
(A) 2
(B) 2
(C) (D) 0
|z + 1| |z + 1| |z + 1|2
Answer and Comments: (D). When we want to find the real (or
the imaginary) part of a complex number which is expressed as a ratio
of two complex numbers, say uv , the standard thing is to multiply both
the numerator and the denominator by the complex conjugate of v
uv̄
and thereby rewrite the ratio as for then we have to worry only
|v|2
about the real and imaginary parts of the product uv̄ which are easier
z−1
to find. Doing so in the present case, we first write w = =
z+1

16
 (z − 1)(z̄ + 1)
. To find the real part of the numerator, we expand it
|z + 1|2
as zz̄ − 1 + z − z̄. Now, zz̄ = |z|2 and 1 are real while z − z̄ is always
purely imaginary. Therefore the real part of the numerator is simply
|z|2 − 1 which is 0 since it is given that |z| = 1. So the real part
of w is also 0. The same argument can be given in a polar form by
writing z as eiθ for some θ, which is possible since |z| = 1. In that case
cos θ − 1 + i sin θ (cos θ − 1 + i sin θ)(cos θ + 1 − i sin θ)
w = = . But
cos θ + 1 + i sin θ (cos θ + 1)2 + sin2 θ
the real part of the numerator now is cos2 θ − 1 + sin2 θ which is 0.
There is also a beautiful y
geometric way to do the prob-
lem. In the Argand diagram,
Pz
the complex numbers z, 1 and
−1 all lie on the unit circle.
The numbers −1 and 1 rep- −1 1
O A x
resent the ends of a diameter. B

In the degenerate case where

z = 1, w and hence its real
part are 0. Otherwise, −1, z
and 1 are vertices of a right an-
gled triangle. Now, the com-
plex numbers z − 1 and z + 1
represent the directed sides AP
and BP of the triangle AP B.
As they are perpendicular to
each other, by Exercise (2.24),
their ratio is purely imaginary.
Hence its real part is 0.
The irony of the problem is that a student can score on it without
doing any of this work. All he has to do is to note that for z = 1,
none of (A), (B), (C) are correct. Hence if at all one of the answers
is correct, it has to be (D)! Hands down. The papersetters could have
|z − 1|
prevented this unwarranted sneak path by giving as one of the
|z + 1|2
false alternatives.

Q. 19 The orthocentre of the triangle with vertices (0, 0), (4, 0) and (3, 4) is

17
 (A) (5/4, −4/3) (B) (3, 12) (C) (3, 5/4) (D) (3, 3/4)
Answer and Comments: (D). Those who have an uncanny memory
can use the formula given in Comment No. 3 of Chapter 8 which
explicitly gives the coordinates of the orthocentre of a triangle in terms
of those of its vertices. The rest of the mortals can proceed directly. The
orthocentre is the intersection of any two altitudes of the triangle. Call
the given vertices as O, A, B respectively. The equation of the altitude
through O can be written down just by inspection since its slope is the
negative reciprocal of the slope of AB which is −4. Hence the equation
of the altitude through O is y = 14 x. Another simple feature of the
problem is that O and A lie on the x-axis. So the altitude through B
is parallel to the y-axis and hence has equation x = 3. Solving these
two equations simultaneously gives (3, 34 ) as the orthocentre.
Although the method is straightforward, the problem is an excel-
lent example of how the calculations can be simplified considerably by
taking advantage of the special features of the problem. (In the present
case there were two such simplifying features, viz., one vertex was at
the origin and one side of the triangle was along the x-axis.) Regret-
tably, the alternatives given also permit an unintended sneak solution.
Once we know that the orthocentre lies on the line y = 14 x we see in-
stantaneously, that out of the four given choices only one, viz., (3, 3/4)
has this property. So, without doing anything further it is the right
choice. It would have been better if one of the wrong choices included
something like (2, 1/2) which also lies on the line y = 14 x, just as the
false answers (3, 12) and (3, 5/4) have been included to make things
difficult for those who go only by the fact that the orthocentre lies on
the line x = 3.

π
 √ tan2 α
Q. 20 For every α ∈ 0, , the value of x2 + x+ √ , x > 0 is greater
2 x2 + x
than or equal to
5
(A) 2 (B) 2 tan α (C) (D) sec α
2
Answer and Comments: (B). It is important to first understand the
problem correctly. α is a fixed real number in the interval (0, π2 ). For
such a fixed α, the given expression is a function of the variable x ∈

18
 (0, ∞) . The problem asks you to check which of the four √ expressions
is a lower bound on the values of this function. If we call x2 + x as
2
u, then the given expression is u + Ku where K = tan α is a positive
constant. Here
q u is also positive. By the A.M. - G.M. inequality,
K2 2
u + u ≥ 2 u × Ku = 2K = 2 tan α.
As a multiple choice question, the answer ends here. But for the
sake of completeness we should also show that none of the other three
alternatives is a lower bound for the given function of x. This follows
2
becasue, in the A.M. - G.M. inequality, equality does hold when u = Ku ,
i.e. when u = K, i.e. when x2 + x = tan2 α. As x varies over (0, ∞)
f (x) varies continuously over (0, ∞). So, by the Intermediate Value
Property for continuous functions, for each given α ∈ (0, π2 ), there do
exist values of x for which the given expression indeed equals 2 tan α.
But if α is sufficiently close to 0, then 2 tan α is less than all the three
expressions 2, 25 and sec α. Hence none of these three expressions can
be a lower bound for the given function for all positive x.
√ tan2 α
The peculiar form of the expression as x2 + x + √ , and
x2 + x
the presence of tan2 α in it coupled with that of 2 tan α in one of the
alternatives make it rather obvious that in this problem the key idea
is the A.M. - G.M. inequality. And like most other problems based on
this inequality, once it strikes you, there is little else left in the problem.
x2 √
Q. 21 A tangent is drawn to the ellipse + y 2 = 1, at the point (3 3 cos θ, sin θ),
27
where 0 < θ < π/2. The sum of the intercepts of the tangent with the
coordinate axes is least when θ equals
π π π π
(A) (B) (C) (D)
6 3 8 4
Answer and Comments: (A). This is a straighforward problem and
the best strategy is to actually find the two intercepts and minimise
their sum as a function of θ over the interval (0, π2 ). The equation of the
√
(3 3 cos θ)x (sin θ)y
tangent at the given point on the ellipse is + = 1.
27 1
As our interest is in the intercepts with the axes, we recast this into
x y
the intercepts form as √ + = 1. Hence the sum of
√ 3 3 sec θ cosec θ
the intercepts is 3 3 sec θ + cosec θ = f (θ) (say). Then f ′ (θ) =

19
 √ √
3 3 sec θ tan θ − cosec θ cot θ. f ′ (θ) = 0 gives 3 3 sin3 θ = cos3 θ, or
tan θ = √13 which means θ = π6 since θ ∈ (0, π2 ). Again, this is sufficient
work in a multiple choice question. But otherwise a rigorous proof that
at this value f atrains its minimum has to be given by showing that
f ′ (θ) < 0 for 0 < θ < π6 and f ′ (θ) > 0 for π6 < θ < π2 .

Q. 22 Let (0, 0), (21, 0) and (0, 21) be the vertices of a triangle. The number
of points having integer coordinates which are strictly inside the given
triangle is

(A) 231 (B) 105 (C) 190 (D) 133

Answer and Comments: (C). The key idea is to see (with a little
help from a diagram if necessary) that the answer is the number of
ordered pairs of positive integers (x, y) for which x + y ≤ 20. Let
k be an integer with 1 ≤ k ≤ 20. Then for x = k, the possible
values of y are from 1 to 20 − k. Hence the desired number is the sum
20 19
(20 − k). Calling 20 − k as r, this is the same as r which equals
P P
k=1 r=0
1
2
× 19 × 20 = 190.

Q. 23 The area of the quadrilateral formed by the tangents to the ellipse

x2 y 2
+ = 1, at the ends of each of its latus-rectum, is
9 5
27 27
(A) (B) 9 (C) (D) 27
4 2
Answer and Comments: (D). The question is straightforward ex-
cept that it involves a relatively peripheral concept, viz., that of a latus
rectum of an ellipse. This is a line through a focus which is perpendic-
x2 y 2
ular to the major axis. For an ellipse in the standard form 2 + 2 = 1
a b
(with 0 < b < a), the foci are at (±ae, 0) where e, the eccentricity,√is
given by b2 = a2 (1 − e2 ). In the present problem, a = 3 and b = 5
which gives e = 23 and hence the foci are at (±2, 0). The latus recta
meet the ellipse at the four points (±2, ± 53 ). By symmetry, the desired
quadrilateral is a rhombus with diagonals along the axes. Hence its
area is four times that of a triangle whose sides lie along the axes and
2x (5/3)y
one of these four tangents. The tangent at (2, 53 ) is + = 1.
9 5
As we are interested in the area of the triangle with two sides as the

20
 intercepts by the axes, we recast this equation into the intercepts form
x y
as + = 1. Hence the desired area is 4 × 12 × 92 × 3 = 27.
9/2 3
Q. 24 The function
Z x2 +1 2
f (x) = e−t dt, x ∈ (−∞, ∞)
x2

is increasing

(A) in (0, ∞) (B) in (−∞, 0) (C) in (−2, 2) (D) nowhere

Answer and Comments: (B). The key idea here is to find the
derivative f ′ (x) and then decide on which of the given intervals its
sign is positive. To find f ′ (x) we apply Formula (19) of Chapter 17
1 1
 
′ −(x2 +1)2 −x4
and get f (x) = e 2x − e 2x = 2x (x2 +1)2 − (x2 )2 . Since
e e
2 2
x + 1 > x , the expression in the parentheses is negative for all x. So
f ′ (x) will be positive if and only if x is negative, i.e. x ∈ (−∞, 0).
π
r
Q. 25 The natural domain of the function f (x) = sin−1 (2x) + is
6
(A) [−1/4, 1/2] (B) [−1/2, 1/4] (C) [−1/4, 1/4] (D) [−1/2, 1/2]
Answer and Comments: (A). In order for sin−1 (2x) to be defined
we must
q have 2x ∈ [−1, 1] and hence x ∈ [−1/2, 1/2]. Further, in order
that sin−1 (2x) + π6 be defined, it is necessary that sin−1 (2x) ≥ − π6
which is the case if and only if 2x ≥ sin(− π6 ), i.e. 2x ≥ − 12 . The two
restrictions are simultaneously satisfied if and only if x ∈ [−1/4, 1/2].

Q. 26 In the interval [0, 1], the mean value theorem is NOT applicable to the
function

1 1


 − x, x<
 2 2

(A) f (x) = 1
2
1
−x , x ≥



2 2
sin x


, x 6= 0
(B) f (x) = x
1, x=0


21
 (C) f (x) = x|x|
(D) f (x) = |x|
Answer and Comments: (A). In order to apply the mean value the-
orem to a function f (x) defined on the interval [0, 1], it is necessary
that the function f (x) be continuous on the closed interval [0, 1] and
differentiable on the open interval (0, 1). So the present problem is not
so much a question on the mean value theorem as on testing the con-
tinuity and the differentiability of the given functions. It is easily seen
that all the four functions are continuous on [0, 1]. So in this respect
they all behave nicely. Note that on (0, 1), the functions in (C) and
(D) are the same as x2 and x respectively. Hence both are differen-
tiable on (0, 1). (The function |x| is not differentiable at 0, but that is
irrelevant in the present problem. Such functions are often included to
mislead a student who vaguely remembers that they fail to be differen-
tiable somewhere.) As for (B) too, sin x and x and hence their ratio are
differentiable on (0, 1). Hence by elimination, (A) is the right choice.
But again, for an honest answer one must indeed show that the func-
tion f (x) is not differentiable at at least one point of (0, 1). Obviously,
the right candidate is the point 12 . The left handed derivative at 12 is
1
− ( 21 + h) − 0
lim− 2 which equals −1 while the righthanded derivative
h→0 h
( 1 − 1 − h)2 − 0
at 12 is lim+ 2 2 = lim+ h which equals 0. Hence f (x) is
h→0 h h→0
not differentiable at 21 . Therefore the mean value theorem cannot be
applied. (In such questions, where the only way to reach the answer
is by testing the given alternatives one after the other, there is an el-
ement of luck depending on the order in which you try the options.
For example, if you notice the non-differentiability of the function in
(A) at 12 , you need not even read the functions that are given in the
other options. Like a good detective, a student should try the most
promising option first.)

Q. 27 Two numbers are drawn at random, one after another and without re-
placement, from the set {1, 2, 3, 4, 5, 6}. The probability that minimum
of the chosen numbers is smaller than 4 is
1 14 1 4
(A) (B) (C) (D)
15 15 5 5

22
 Answer and Comments: (D). The key idea is to find the comple-
mentary probability, which is easier. There are 6 × 5 = 30 possible
ways of drawing the numbers. The given event will fail to occur when
both the numbers are from the subset {4, 5, 6} and this can happen in
6
3 × 2 = 6 ways. Hence the complementary probability is 30 = 15 . So
the desired probability is 1 − 15 = 45 .

x2 y2 π
 
Q. 28 For the hyperbolas given by 2
− 2 = 1, α ∈ 0, , which of
cos α sin α 2
the following remains unchanged with varying α?

(A) eccentricity (B) abscissae of the foci

(C) equations of the directrices (D) abscissae of the vertices

Answer and Comments: (B). Another question which seems to be

designed to test the knowledge of definitions of terms some of which
are rarely used today. (For example, ‘abscissa’ simply means the x-
coordinate and ‘abscissae’ is its chaste plural!) Now for the hyperbola in
x2 y 2
the standard form 2 − 2 = 1, the foci are at the points (±ae, 0) and
a b
the equations of the directrices are x = ± ae where e, the eccentricity,
is given by b2 = a2 (e2 − 1). The key idea is therefore, to first find
the eccentricity of a typical member of the given family of hyperbolas.
Taking a = cos α and b = sin α gives e2 = tan2 α + 1 = sec2 α. So
e = sec α. Thus the eccentricity does change with α. But the abscissae
of the foci are ± cos α sec α i.e., ±1. These do not depend upon α. So
(B) is the correct answer and we need not scan further. In particular,
we are spared the consequences of not knowing what vertices mean!
In a limited time examination, it would, in fact, be foolish to spend
time on such redundant options once the correct alternative is found.
Still, for the sake of knowledge, the equations of the directrices are
x = ± cos2 α which depend on α. As for vertices, they are the points
where the hyperbola cuts the major axis, i.e. the points (± cos α, 0) in
the present case. They too depend on α.

23
 MAIN PAPER OF JEE 2003

Problem 1. Let z1 , z2 be complex numbers such that |z1 | < 1 < |z2 |. Show
that
 1 − z1 z2 
 
 <1

z1 − z2 
Analysis and Solution: Geometrically, the data means that z1 lies inside
while z2 lies outside the unit circle. This in particular implies that they
1 − z1 z2
cannot be equal and so the complex number is certainly defined. We
z1 − z2
have to show that this number also lies inside the unit circle. Unfortunately
there is no obvious geometric interpretation to this number. And so instead
of tackling the problem geometrically, let us tackle it algebraically using
standard properties of complex numbers and inequalities. As is the case
with many other problems with inequalities, the success may come if we go
on replacing the desired inequality with equivalent (or stronger) inequalities
 1 − z1 z2 

and prove the last one. In the present problem, we first use that   =
z1 − z2 
|1 − z1 z2 |
. Hence the problem is reduced to proving that
|z1 − z2 |

|1 − z1 z2 | < |z1 − z2 | (1)

There is no standard inequality which would imply this and so we have

to start from scratch. Since both the sides are non-negative, it is a good
idea to square them, because the square of the absolute value of a complex
number is amenable to algebraic manipulations. So the problem is reduced
to showing

|1 − z1 z2 |2 < |z1 − z2 |2 (2)

We now use the fact that for any complex number z, |z|2 is simply zz.
Using standard properties of complex conjugation, the problem now reduces

24
to proving

(1 − z1 z2 )(1 − z1 z2 ) < (z1 − z2 )(z1 − z2 ) (3)

If we expand both the sides and cancel the common terms, this reduces
to showing

1 + |z1 |1 |z2 |2 < |z1 |2 + |z2 |2 (4)

or equivalently to showing that

0 < |z1 |2 + |z2 |2 − |z1 |2 |z2 |2 − 1 (5)

But the R.H.S. factors as (1 − |z1 |2 )(|z2 |2 − 1) which is positive since the
hypothesis implies that both the factors are positive. Hence all inequalities
we have listed hold true including the one to be proved.

Problem 2. Let ai , i = 1, 2, . . . , n, be complex numbers with |ai | < 2 for

each i. Prove that there is no complex number such that
n
1
ai z i = 1
X
|z| < and
3 i=1

Analysis and Solution: The phrasing of the question is somewhat defec-

tive. Instead of saying ‘there is no complex number such that’ it should have
said ‘there is no complex number z such that’. Also since complex numbers
are involved, the use of i as an index of summation is not a good idea, as it
can be confused with the complex number i.
It is tempting to consider the polynomial p(z) = an z n + an−1 z n−1 + . . . +
a1 z − 1. The goal then is to show that this polynomial has no root in the
disc {z : |z| < 31 }. Unfortunately, even though the Fundamental Theorem
of Algebra guarantees that a complex polynomial of degree n has n complex
roots (counted according to multiplicities), it says virtually nothing about
the location of these roots. So this approach is not likely to work.
Since the only information we are given is in the form of upper bounds
on the absolute values of the numbers a1 , a2 , . . . , an and that of |z|, the
only way to tackle the problem is through inequalities involving absolute
values of complex numbers. The foremost such inequality is the triangle

25
inequality which asserts, in a generalised form, that for any complex numbers
z1 , z2 , . . . , zn ,

|z1 + z2 + . . . + zn | ≤ |z1 | + |z2 | + . . . + |zn | (1)

Let us see how this inequality comes into the picture. Assume there is a
complex number z such that
n
1
ai z i = 1
X
|z| < and (2)
3 i=1

holds true. Then in particular, taking absolute values we must also have
n
 
X 

 ai z i  =1 (3)
 
i=1

If we apply (1) to the L.H.S. of (3), we get

1 ≤ |a1 z| + |a2 z 2 | + . . . + |an z n | (4)

which can be further rewritten (using the fact that the absolute value of a
product is the product of the absolute values) as

1 ≤ |a1 ||z| + |a2 ||z|2 + |a3 ||z|3 + . . . + |an ||z|n (5)

At first sight it does not appear that we can get any contradiction from
this because the integer n can be quite large. But the saving feature is that
since |z| < 1, the higher powers of |z| will be very small in magnitude and
hence even though the number of terms is large, the total contribution may
not be that large. Once this idea strikes, a careful estimation of the terms on
the R.H.S. of (5) leads to a solution. We are given that |ak | < 2 for every k.
Further we also have |z| < 13 and hence |z|k < ( 31 )k for every positive integer
k. Hence |ak ||z k | < 2( 13 )k for k = 1, 2, . . . , n. Using this in (5) gives

1 1 1
1 ≤ 2( + ( )2 + . . . + ( )n ) (6)
3 3 3
Using the formula for the sum of a geometric progression gives

1 1 − ( 31 )n 1
1≤2× × 1 = 1 − ( )n (7)
3 1− 3 3

26
which is a contadiction since ( 13 )n > 0.
The problem is an excellent combination of inequalities and geometric
progression. Once the key idea strikes, there is not much numerical work.

Problem 3: Let a, b, c be positive real numbers with abc = 1. Let

 
a b c
 b c a 
A= 

c a b

If AT A = I, where AT is the transpose of A and I is the identity matrix,

then determine the value of a3 + b3 + c3 .

Analysis and Solution: The reference to the transpose AT is redundant,

because since A is symmetric, AT equals A and we could as well have been
given that A2 = I. A direct computation gives
a2 + b2 + c2 ab + bc + ca ab + bc + ca
 
T 2
A A = A =  ab + bc + ca a2 + b2 + c2 ab + bc + ca  (1)
 

ab + bc + ca ab + bc + ca a2 + b2 + c2
Equating this with I, the identity matrix, gives a system of two equations
a2 + b2 + c2 = 1 (2)
and ab + bc + ca = 0 (3)
But since a, b, c are given to be positive, (3) can never hold. Hence the data
of the problem is inconsistent. Let us, however, ignore this inconsistency by
dropping the hypothesis that a, b, c are all positive. In that case, (2) and (3)
together imply that (a + b + c)2 = a2 + b2 + c2 = 1. Therefore we have
a + b + c = ±1 (4)
(If we assume that a, b, c are positive then we must discard the negative sign.
However, as just noted such an assumption is not consistent with (3).)
We are given that abc = 1. We have thus got hold of the numerical values
of the three elementary symmetric functions of a, b, c. We have to find the
value of a3 + b3 + c3 which is also a symmetric function of a, b, c. To do this,
we use the well-known factorisation
a3 + b3 + c3 − 3abc = (a + b + c)(a2 + b2 + c2 − ab − bc − ca) (5)

27
From (2), (3), (4) and abc = 1, we thus get a3 + b3 + c3 = 3 ± 1, i.e. 4 or
2. (Again, if we continue to assume that a, b, c are positive, then we get only
one answer, viz., a3 + b3 + c3 = 4.)
The solution officially given on the internet along with the problem con-
sists of first showing that det(A) = ±1 by using the fact that det(AT A) =
(det(A))2 = det(I) = 1. By a direct computation, the determinant of A
comes out to be 3abc − (a3 + b3 + c3 ) which is factorised using (5). The fact
that the last factor in the R.H.S. of (5) is always non-negative, coupled with
the positivity of a, b, c then implies that det(A) = −1. Since abc = 1, we
then get a3 + b3 + c3 as 4.
From this solution it appears that the problem was designed to test the
knowledge of elementary properties of determinants, of the factorisation (5)
and of inequalities (needed to show that the second factor of (5) is non-
negative). However, the whole problem is hypothetical. Even if we drop the
assumption that a, b, c are positive, it can be shown that there are no real
numbers which satisfy (2), (3) and (4). For, if there were, then (depending
upon which sign holds in (4)), they would have to be the roots of the cubic
x3 − x2 − 1 = 0 or of the cubic x3 + x2 − 1 = 0. But it is easy to show (using
a little calculus) that both these cubics have only one real root each. So if
the matrix A given in the problem has the property that AT A = I, then its
entries will have to be complex.
It is instructive to paraphrase the condition AT A = I slightly. Matrices
(with real entries) which have this property are very important and are called
orhogonal. The peculiar name comes from the concept of orthogonality, i.e.
perpendicularity of vectors. To see this, suppose A = (aij ) is a 3 × 3 (real)
matrix with the property that AT A = I. Writing this out fully, we get
    
a11 a21 a31 a11 a12 a13 1 0 0
AT A =  a12 a22 a32   a21 a22 a23  =  0 1 0  (6)
    

a13 a23 a33 a31 a32 a33 0 0 1

Now let us identify the columns of A with three dimensional vactors. Thus,
a11
for example, the first column of A, viz.,  a21  will be identified with the
 

a31
vector a11 i + a21 j + a31 k. Let us denote this vector by v1 . Similarly, identify
the second and the third column of A with vectors v2 and v3 .
Now if we carry out the multiplication of the two matrices in the middle
term of (6), we see that the entry in the (i, j)-th place will be simply the dot

28
product of the vectors vi and vj . For example, the entry in the second row
and the first column of the product matrix is a12 a11 + a22 a21 + a32 a31 and this
is precisely the dot product v2 · v1 . Doing this for all entries, we can rewrite
(6) as
   
v1 · v1 v1 · v2 v1 · v3 1 0 0
AT A =  v
 2 · v1 v2 · v2 v2 · v 
3  =  0 1 0 
 
(7)
v3 · v1 v3 · v2 v3 · v3 0 0 1
If we equate the corresponding entries of the matrices we have got, we see
that (7) simply says that the dot product vi ·vj equals 1 if i = j and 0 if i 6= j.
In other words, every two distinct columns of the matrix A are orthogonal
to each other when treated as vectors. This explains why such a matrix is
called orthogonal. Another way to say that a matrix is orthogonal is to say
that its column vectors form an orthonormal basis. As these concepts can be
generalised for higher dimensional vectors too, we see that this paraphrase
is valid for every n, and not just for n = 3.
In the light of this new interpretation, the inconsistency of the data be-
comes even more apparent. For, suppose A is a 3 × 3 matrix such that
AT A = I. To say that the entries of A are all positive is equivalent to saying
that all its column vectors lie in the first orthant. (See Comment No. 10
of Chapter 21 for a definition.) But no two vectors in the same orthant can
be orthogonal to each other unless they lie in the boundary planes of the
orthant (in which case some of their entries will have to be 0).

Problem 4: Find the point on the ellipse x2 + 2y 2 = 6 closest to the line

x + y = 7.

Analysis and Solution: There are several ways of doing this problem.
x2 y 2
We first rewrite the ellipse in the standard form as + = 1. Hence
√6 3 √
a typical point, say P , on the ellipse is of the form ( 6 cos θ, 3 sin θ) for
some θ √∈ [0, 2π]. The distance of this point from the line x + y = 7 is
√
| 6 cos θ+√ 3 sin θ−7|
√
2
. Minimising this is equivalent to minimising ( 6 cos θ +
√
3 sin θ − 7)2 as a function of θ. Call this function as f (θ). Then
√ √ √ √
f ′ (θ) = 2( 6 cos θ + 3 sin θ − 7)( 3 cos θ − 6 sin θ) (1)
√ √
√Note√that the first factor can never vanish since | 6 cos θ + 3 sin θ| ≤
( 6 + 3) < 7 for every θ. (Or, equivalently, because the line x + y =

29
7 does not intersect the given ellipse.) The second factor vanishes when
tan θ = √12 . This gives two possible values of θ in the interval [0, 2π]. The
two corresponding points on the ellipse represent the closest and the farthest
points from the given line. From tan θ = √12 , we have sin θ = ± √13 and
√
cos θ = ± √23 . The corresponding points on the ellipse are (2, 1) and (−2, −1).
Clearly, (2, 1) is closer to the line x + y = 7. So this is the point on the given
ellipse which is closest to the line x + y = 7.
The essence of the solution above was to represent the points of the ellipse
in terms of a single parameter θ and thereby reduce the problem to that of
minimising a certain function of one variable, viz., θ. However, as mentioned
at the end of Comment No. 14 of Chapter 13, instead of taking the parametric
equations of the ellipse, we can leave it in the form x2 +2y 2 = 6 and apply the
method of Lagrange’s multipliers. Although this method is not in the JEE
syllabus (since it requires the differentiation of functions of two variables),
for those who are familiar with it, the problem amounts to minimising the
function 12 (x + y − 7)2 (which represents the square of the distance of a
point (x, y) from the line L) subject to the constraint g(x, y) = 0 where
g(x, y) = x2 + 2y 2 − 6. Taking partial derivatives w.r.t. x and y, we get a
system of equations
1
(x + y − 7) = 2λx (2)
2
1
and (x + y − 7) = 4λy (3)
2
Eliminating the multiplier λ gives x = 2y. Solving this simultaneously with
x2 + 2y 2 = 6 gives the same points, viz. (2, 1) and (−2, −1) as before. The
first point is the closest and the second the farthest from L.
There is also a geometric way to find the answer analogous to the alternate
solution to the problem in Comment No. 14 of Chapter 13. Suppose P is a
point on the ellipse which is closest to the line (say L) with equation x+y = 7.

30
Let L′ be the line through y
P that is parallel to L. We
claim that all points of the
L
ellipse lie only on one side P
of L′ . For othrwise, there
will be points on the ellipse x
O
which are closer to the line
L’
L than P is. Put differ-
ently, this means that L′
must touch the ellipse at
P . (A similar statement also
holds for the point on the el-
lipse which is farthest from
L.)
As a result, to find P , we look for a point on the ellipse where the tangent
is parallel to L, i.e., has slope −1. Letting P = (h, k), the equation of the
tangent at P is hx 6
+ ky
3
= 1. Hence its slope is − 2k h
. (This can also be
2 2
obtained by differentiating the equation x + 2y = 6 implicitly w.r.t. x.) So
if (h, k) is the point on the ellipse which is closest to L, we have

h = 2k (4)

Also, since (h, k) lies on the ellipse we have

h2 + 2y 2 = 6 (5)

Solving (2) and (3) simultaneously, we get h = ±2 and k = ±1. This gives
two points, viz., (2, 1) and (−2, −1) on the ellopse. The first one is closest
while the second one is farthest from the line L.

Problem 5. A person takes three tests in succession. The probability of

his passing the first test is p, that of his passing each successive test is p or
p/2 according as he passes or fails in the preceding one. He gets selected
provided he passes at least two tests. Determine the probability that the
person is selected.

Analysis and Solution: We might think of each test as a toss of a coin and
identify passing the test with a head showing and failure with a tail. In that
case, the problem amounts to finding the probability of getting at least two

31
heads in three tosses. Such prblems were studied under binomial  probability
 
3 3
in Comment No. 13 of Chapter 22. The answer would be 3 p + 32 p2 (1 −
p) = 3p2 − 2p3 .
But this is a wrong answer. There is a vital difference between the given
problem and the problem of getting two heads in three tosses. In the latter,
the outcome of each toss is independent of that of the previous one. In the
problem asked, on the other hand, the probability of passing a test depends
on the outcome of the previous test. So, the trials are not independent.
Therefore binomial probability is not applicable. (This is a point to keep in
mind. Although it is definitely an asset to be able to recall similar problems,
one should not be misled by superficial similarity. Even one vital difference
can completely change the problem.)
So we have to think fresh. Let P and F denote passage and failure in a
given test. As there are three tests, there are 8 possible outcomes, denoted
by ordered triples like P P F or F P F etc. The probability of each can be
found from the data. For example, P r(P P F ) = p × p × (1 − p). But
P r(P F P ) = p × (1 − p) × p2 and so on. We are interested in those outcomes
where at least two P ’s occur. They are : P P P , P P F , P F P and F P P . They
are mutually exclusive. So the desired probability is simply the sum of the
probabilities of these four outcomes. It comes out as p3 + p2 (1 − p) + p(1 −
p) 2p + (1 − p) p2 p = 2p2 − p3 .
Normally, probability problems involve independent trials. The present
problem is different and innovative. In fact, it is an excellent example where
going blindly by similarity will fail but cool, fresh thinking will win.
It is natural to inquire what the answer will be if instead of 3 tests, we
had n tests. We can also change the selection criterion to require at least
m passings instead of 2. The case m = 1 is easily handled by taking the
complementary probability (all failures) which is (1 − p)(1 − 2p )n−1 . Hence
the desired probability in this case is 1 − (1 − p)(1 − p/2)n−1 .
The case m = 2 is fairly complicated for a general n. Again we take the
complementary probability, i.e. the probability of the person not getting
selected. This can happen when either he fails all tests or passes exactly one.
The probability of the former was already computed above as (1 − p)(1 −
p/2)n−1. For brevity denote this by fn . For k = 1, 2, . . . , n, let Ek be the
event that the person passes only the k-th test and fails all others. Then the

32
complementary probability we are after is the sum
n
X
fn + P (Ek ) (1)
k=1

We already computed fn above. The computation of P (Ek ) is slightly differ-

ent for k = 1 and n than for 1 < k < n. Using the same reasoning as for the
computation of all failures, it is clear that

P (E1 ) = pfn−1 = p(1 − p)(1 − p/2)n−2 (2)

and P (En ) = fn−1 (p/2) = (1 − p)(1 − p/2)n−2(p/2) (3)

However, for 2 ≤ k ≤ n − 1, we have

P (Ek ) = fk−1 (p/2)fn−k = (p/2)(1 − p)2 (1 − p/2)n−3 (4)

Note that this is independent of k. So among the terms under the sigma
sign in (1), all the terms except the first and the last equal (p/2)(1 − p)2 (1 −
p/2)n−3 each. Also the first term is twice the last one. We thus get that the
probability of at most one success in n tests as

(1 − p)(1 − p/2)n−1 + (3p/2)(1 − p)(1 − p/2)n−2

+ (n − 2)(p/2)(1 − p)2 (1 − p/2)n−3 (5)

which simplifies to
(1 − p)(1 − p/2)n−3
(2 + (n − 1)p(1 − p)) (6)
2
The next case m = 3 is far too comolicated to be given here. However,
the special case of (6) obtained by taking n = 3 gives (1−p) 2
(2 + 2p − 2p2 )
2 2 3
= (1 − p)(1 + p − p ) = 1 − 2p + p . This is the probability of passing at most
one test in 3 tests. Therefore subtracting it from 1 we get the probability of
getting selected as 2p2 − p3 . This gives an alternate solution to the problem
asked. Of course, nobody would recommend this solution in an examination.
But in your leisure time it is always a good idea to try to get the answer by
other methods, some of which are more general than what is needed.

Problem 6. In a combat, A targets B, and both B and C target A. The

2 1 1
probabilities of A, B, C hitting their targets are , , and respectively.
3 2 3
33
They shoot simultaneously and A is hit. Find the probability that B hits his
target whereas C does not.

Analysis and Solution: The first thing to realise is that the probability
of A’s hitting his target is completely irrelevant in this problem. (This is a
point to note. Normally, every piece of the data is needed in the solution.
But once in a while, the papersetters can test your ability to pick up only
the relevant things from the data.)
Let E and F be, respectively, the events that B and C hit A. We are given
that P (E) = 21 and P (F ) = 13 . We are further given that at least one of E
and F occurs. The problem amounts to finding the conditional probability
(E ∩ F ′ |E ∪ F ), i.e. the probability that E ∩ F ′ occurs, given that E ∪ F
occurs. (Here the dash ′ denotes the complementary event.) To find this we
can apply the law of conditional probability (Comment No. 5, Chapter 22)
to get

′ P (E ∩ F ′ ∩ (E ∪ F )) P (E ∩ F ′ )
P (E ∩ F |E ∪ F ) = = (1)
P (E ∪ F ) P (E ∪ F )

where the second equality follows from the fact that E ∩ F ′ is a sub-event of
E ∪ F . We are given P (E) and P (F ). Since E and F are independent, so
are E and F ′ . Hence

1 1 1 2 1
P (E ∩ F ′ ) = P (E)P (F ′) = (1 − ) = × = (2)
2 3 2 3 3
Finally, to get the denominator, viz. P (E ∪ F ), we consider the comple-
mentary probability P ((E ∪ F )′ ) which is the same as P (E ′ ∩ F ′ ). Since E,F
are mutually independent, so are E ′ and F ′ . So, P (E ′ ∩ F ′ ) = P (E ′ )P (F ′) =
(1 − 12 )(1 − 31 ) = 31 . Hence,

1 2
P (E ∪ F ) = 1 − P ((E ∪ F )′ ) = 1 − = (3)
3 3
1/3 1
= .
Substituting (2) and (3) into (1) gives the desired probability as
2/3 2
The problem is a good combination of complementary probability, condi-
tional probability and independence of events

34
Problem 7. Let f : [−2a, 2a] → IR be an odd function such that the left
derivative of f at x = a is zero. If f (x) = f (2a − x) for x ∈ (a, 2a), then
compute the left derivative of f at x = −a.

Analysis and Solution: Here the function f satisfies two functional equa-
tions. On the entire interval [−2a, 2a], it satisfies the condition f (−x) =
−f (x). But on the subinterval (0, 2a) it satisfies the relation that f (x) =
f (2a − x), which is analogous to the definition of an even function. These
conditions have obvious geometric interpretation. The first condition implies
(and is equivalent to saying) that the graph of the function f (x) is symmet-
ric about the origin. Consequently, as pointed out in Exercise (18.6)(d), the
tangent at a point (x0 , f (x0 )) is parallel to that at the point (−x0 , f (−x0 )).
Hence f ′ (−x0 ) = f ′ (x0 ). More generally, even if we are given only the left
and the right derivatives at x0 , then by the symmetry of the curve y = f (x)
w.r.t. the origin, we get,

f−′ (−x0 ) = f+′ (x0 ) (1)

and f+′ (−x0 ) = f−′ (x0 ) (2)

(See figure (a) next page.)

y

− x0 O
. . x . . .
x0 2 a − x0 a x0

(a) (b)

Similarly, the second condition viz. f (2a − x) = f (x) for all x ∈ (a, 2a)
(which continues to hold trivailly even for x = a) geometrically means that
the graph is symmetric about the line x = a as shown in Figure (b) above.
Hence the tangents at the points (x0 , f (x0 )) and (2a − x0 , f (2a − x0 )) have
opposite slopes. And more generally,

f+′ (x0 ) = −f−′ (2a − x0 ) (3)

and f−′ (x0 ) = −f+′ (2a − x0 ) (4)

35
 With these observations, f−′ (−a), i.e. the left handed derivative of at a
can be computed as follows. Taking x0 = a in (1), we get f−′ (−a) = f+′ (a).
But by putting x0 = a in (3), f+′ (a) = f−′ (a) which is given to be 0. So
f−′ (−a) is 0.
This solution was based on the geometric interpretations of the given
conditions and those of the left and right derivatives. An analytical proof
essentailly amounts to proving the properties (1) and (3) above. Thus we are
given that
f (a) − f (a − h)
lim =0 (5)
h→0+ h
Now,
f (−a) − f (−a − h)
f−′ (−a) = lim+
h→0 h
−f (a) + f (a + h)
= lim+ , since f is odd
h→0 h
−f (a) + f (2a − a − h)
= lim+ , since f (2a − x) = f (x)
h→0 h
−f (a) + f (a − h)
= lim+ = 0 using (5)
h→0 h
The analytical solution is short and logical than the geometric solution.
In fact, as commented after Theorem 3 in Chapter 13, an argument based on
geometric intuition can never be a substitute for a rigorous analytical proof,
no matter how convincing it is to a layman. Still, the fact remains that one
often thinks of the right ideas in terms of geometric interpretation and then
gives it an analytical expression. Without the aid of the geometric interpre-
tation, it is doubtful if anybody can think of the manipulations needed in
the analytical solution.

Problem 8: Let p(x) be a polynomial such that p(1) = 0 and

d
(p(x)) > p(x) for all x > 1.
dx
Show that p(x) > 0 for all x > 1.

Analysis and Solution: This is a problem which combines differential

equations with inequalities. Since p(1) is given to be 0, in order to prove the

36
assertion, it would suffice to show that p′ (x) > 0 for all x > 0, for in that
Rx
case one can write p(x) = p(x) − p(0) = p′ (t) dt > 0 as the integrand is
0
positive on the interval [0, x] except at the end point 0. Now, we are given
that p′ (x) > p(x) for x > 1 and so in order to show that p′ (x) > 0 for x > 1,
it would suffice to show that p(x) > 0 for x > 1. This is the same as the
assertion we want to prove in the first place. So, this line of attack is no
good because it reduces the problem to some other problem whose solution
hinges on that of the original problem. (In logic, such a situation is called a
vicious cycle.)
p′ (t)
Let us try some alternative. It is tempting to rewrite the data as >0
p(t)
and then integrate both the sides w.r.t. t from 1 to x, keeping in mind that
the L.H.S. is the derivative of ln p(t) w.r.t. t. But the trouble is that p(1) = 0
p′ (t)
and so even if we ignore the endpoint 1, the ratio will, in general, not
p(t)
be a bounded function on the open interval (1, x).
The difficulty encountered here is analogous to that in the problem solved
in Comment No. 13 of Chapter 19. And predictably, the trick to resolve it
is also analogous. If we multiply the given inequality by e−x , we can recast
it as an inequality about derivatives. Indeed, if we call e−x p(x) as f (x), then
f ′ (x) is simply e−x (p′ (x) − p(x)). As the exponential function assumes only
positive values, and the factor (p′ (x) − p(x) is given to be positive for all
x > 1, we get

f ′ (x) > 0 for all x > 0 (1)

Note further that f (0) = p(0) = 0. Hence by integration, for every x > 0 we
get
Z x
f (x) = f (x) − f (0) = f (t) dt > 0 (2)
0

for all x > 0. writing back as f (x) = e−x p(x) and multiplying by e−x (which
is always positive) gives that p(x) > 0 as desired.
The solution is somewhat tricky and only those who are familiar with
similar problems such as the one referred to above are likely to get it. Note
also that no use is made of the hypothesis that p(x) is a polynomial. It could
have been any continuously differentiable function of x. In fact, it would
have been better, had the hypothesis that p(x) is a polynomial been dropped

37
from the statement of the problem, because, as it stands no polynomial p(x)
with p(x) = 1 can have the property that p′ (x) > p(x) for all x > 1. This
can be seen intuitively as follows. The order of growth of the absolute value
of a polynomial function is determined by its degree. The derivative of a
polynomial has lower degree than that of the original polynomial. So the
only way it can be bigger than the original polynomial for all large x is when
the leading coefficient of the original polynomial is negative. But in that case,
the polynomial will assume negative values for all large x, contradicting the
conclusion of the problem. For a formal proof, suppose
p(x) = an xn + an−1 xn−1 + . . . + a2 x2 + a1 x + a0 (3)
where a0 , a1 , . . . , an−1 , an are some real numbers and an 6= 0. Assume further
that p(1) = 0. Then
p′ (x) − p(x) = −an xn + (nan − an−1 )xn−1 + . . . + (a1 − a0 ) (4)
which is a polynomial of degree n too with leading coefficient −an . Now, for
large x, the sign of a polynomial expression depends only on the sign of its
leading coefficient. So if p′ (x) > p(x) for all x > 1, then by (4) we must have
an < 0. But in that case the result of the problem would imply that p(x) > 1
for all x > 1. By (3) this would mean an > 0. This contradiction shows that
the data of the problem is vacuous, like that of Problem 3.
However, there are functions which satisfy the conditions in the problem
except that they are not polynomials. A simple example is the function
2
ex − 1. So if the word ‘polynomial’ in the statement of the problem is
replaced by ‘a continuously differentiable function’, then the problem is not
vacuous.

Problem 9. Let n and k be positive integers such that n ≥ k. Prove that

! ! ! ! ! ! !
n n n n−1 n n−k n
2k − 2k−1 + . . . + (−1)n =
0 k 1 k−1 k 0 k
!
m
where =m C r .
r
  
n n−r
Analysis and Solution: The typical term on the L.H.S. is (−1)r 2k−r r k−r
  
k n n−k
where 0 ≤ r ≤ k. However, then the last term ought to be (−1) k 0
  
n n−k
and not (−1)n k 0
. We assume that this is a misprint and proceed.

38
 Here every term involves the product of two binomial coefficients. There
are many familiar sums where each term involves just one binomial coefficient.
There are very few standard sums in which each term contains a product of
two binomial coefficients. One such identity is given by Equation (12) in
Chapter 5. It says that
k
! ! !
X m n m+n
= (1)
r=0 r k−r k
But in this identity, the upper indices of the binomial coefficients in each
term are fixed (i.e. independent of the term). As this is not the case with
the binomial coefficients in the identity we are asked to prove, it is unlikely
that (1) will be of much use here.
The best bet is, therefore to see if we can reduce the given sum to one in
which each term has only one binomial coefficient. This could be done if we
can extract some binomial coefficient as a common factor from each term.
And the R.H.S. of the identity
  to be proved suggests that this common fac-
tor is most likely to be nk . Once this idea strikes us, we try to rewrite the
    
nn−r
product rk−r
in such a way that nk will be a factor. This is indeed pos-
     
siblebecause nr n−r
k−r
= n k
k r
. (This is precisely the identity in Exercise
(5.1)(i), except for the change of notation. It can be proved algebraically as
well as by a combinatorial argument.)  
The advantage of this rewriting is that now the factor nk is common to
all the terms of the summation and hence can be taken out. So, calling the
given sum as S, we get
k
! !
n X k
S= (−1)r 2k−r (2)
k r=0 r
 
Now (−1)r 2k−r kr is precisely the r-th term in the binomial expansion of
(2 − 1)k (where the term count starts with 0 rather than with 1 and we write
the powers of 2 in decreasing order). Since (2 − 1)k = 1 no matter what k is,
the sum in the R.H.S. of (2) is simply 1. This completes the solution.

Problem 10. If the function f : [−1, 1] → IR is continuous and even, then

show that
Z π/2 √ Z π/4
f (cos 2x) cos xdx = 2 f (sin 2x) cos xdx
0 0

39
 Analysis and Solution: This problem is highly analogous to the last
problem solved in Comment No. 14 of Chapter 18. Such problems are often
based on the identity
Z L Z L
g(x) dx = g(L − x) dx (1)
0 0

π/2
R
Let I = f (cos 2x) cos x dx. Applying (1) to this integral and noting that
0
f (cos 2(π/2 −x)) = f (cos(π −2x)) = f (− cos 2x) = f (cos 2x) since f is given
to be an even function, we get
Z π/2
I= f (cos 2x) sin x dx (2)
0

But it is not immediately clear how to relate this to the R.H.S. because the
latter is an integral over a different interval, viz., [0, π/4] which is the left
half of the interval of integration of I, viz. [0, π/2]. This suggests that we
should perhaps split the given integral I into two parts as I1 + I2 where
Z π/4
I1 = f (cos 2x) cos x dx (3)
0
Z π/2
and I2 = f (cos 2x) cos x dx (4)
π/4

If we apply (1) to I1 we get

Z π/4
I1 = f (cos(π/2 − 2x) cos(π/4 − x) dx
0
1 π/4
Z
= √ f (sin 2x)(cos x + sin x) dx (5)
2 0
√
which is a bit heartening because we are able to involve 2 in it (albeit in
the denominator and not in the numerator as we would like to, but that can
be taken care of if we are able to double the integral).
Let us now work with I2 . We want to convert it to an integral over an
interval beginning at 0. This can be done by a simple substitution u =
x − π/4. Then as x varies over [π/4, π/2], u varies over [0, π/4] and we have
Z π/2
I2 = f (cos 2x) cos x dx
π/4

40
 Z π/4
= f (cos 2(π/4 + u) cos(π/4 + u) du
0
1 π/4
Z
= √ f (− sin 2u)(cos u − sin u) du (6)
2 0

We now use (for the first time) that the function f is even. Therefore we
have f (− sin 2u) = f (sin 2u). Putting this into (6) and replacing the dummy
variable u of integration by x, we get
1 π/4
Z
I2 = √ f (sin 2x)(cos x − sin x) dx (7)
2 0

If we add (5) and (7) some terms add up while some get cancelled and we
are left with,
√ Z π/4
I = I1 + I2 = 2 f (sin 2x) cos x dx (8)
0

which is exactly what we wanted to prove.

The key idea in this proof was to split the integral on the L.H.S. into two
parts as suggested by the form of the R.H.S. Although we abandoned (2), a
π/2
R
solution can also be given using it. We simply add it to I = f (cos 2x) cos x dx
0
and get
Z π/2
2I = f (cos 2x)[sin x + cos x] dx (9)
0
√
We now write sin x + cos x as 2 cos(x − π/4) and get
√ Z π/2
2I = 2 f (cos 2x) cos(x − π/4) dx (10)
0

The next step is to convert this into an integral over the interval [−π/4, π/4]
by the substitution x = u + π/4. The advantage in doing so is that the new
interval is symmetric about the origin and that is especially convenient when
even or odd functions are involved. So, we have
√ Z π/4
2I = 2 f (cos(2u + π/2) cos u du
−π/4
√ Z π/4
= 2 f (− sin 2u) cos u du
−π/4
√ Z π/4
= 2 f (sin 2u) cos u du (11)
−π/4

41
where in the last step we have again used that f is even. Since the sine
function is odd and the function f is even, the composite function f (sin 2u)
is an even function of u. The function cos u is also even. Hence the integrand
of the integral in the R.H.S. of (11) is an even function of u. As the interval of
integration is symmetric about 0, by a well-known property of such integrals
we have
√ Z π/4
2I = 2 2 f (sin 2u) cos u du (12)
0

Cancelling the factor 2 from both the sides and replacing the dummy
variable u by x we get the desired result. It is instructive to compare the two
solutions. Both are based on the idea of splitting the interval of integration
suitably and that of shifting it. In the first solution, we first split the interval
[0, π/2] into two halves and then shifted the right half to the left by π/4 with
a suitable substitution. In the second solution, on the other hand, we first
shifted the interval [0, π/2] to the interval [−π/4, π/4]. Although we then
derived (12) from (11) by a well-known property of integrals, the proof of
this property is again based on the idea of splitting the interval [−π/4, π/4]
into two halves. So, ultimately, the essential ideas behind the two solutions
are the same.

Problem 11. If the numbers a, b, c are in arithmetic progression and a2 , b2 , c2

c
are in harmonic progression, then show that either a = b = c or a, b, − are
2
in geometric progression.

Analysis and Solution: Here we have three variables a, b, c. The given con-
ditions give only two equations involving them. So we cannot determine the
three numbers from them. But as pointed out at the beginning of Comment
No. 6 of Chapter 2, the property of being in any one type of progression is
unchanged if we multiply all terms by any (non-zero) scalar. This allows us
to choose one of the terms arbitrarily. In the present problem, all the three
numbers a, b, c are non-zero (as otherwise the concept of harmonic progres-
sion will not make sense). Without loss of generality, we suppose that b = 1.
(For otherwise, we can work with the numbers ab , 1 and cb which satisfy the
same conditions.)
With this simplification, the data gives

a+c = 2 (1)

42
 and a2 + c2 = 2a2 c2 (2)

We now solve these two equations simultaneously. Squaring the first equa-
tion and using (2) we get

a2 c2 + ac = 2 (3)

We treat this as a quadratic in ac and solve it to get

ac = 1 (4)
or ac = −2 (5)

The first possibility gives c = a1 which along with (1) implies a + a1 = 1. This
can be converted to a quadratic in a, viz., a2 − 2a + 1 = 0 whose only root
is a = 1. But in that case c would also equal 1. Since b is already taken as
1, a, b, c are all equal in this case.
The second possibility, viz., ac = −2 means that a, 1 and − 2c are in a
G.P. But since b = 1, this means a, b, − 2c are in a G.P. This completes the
solution.
This solution was short and elegant because of the trick used in it, viz.
that of taking b as 1. If we cannot think of such a simplifying measure, we
can still manipulate the data in a rather pedestrian way to get the conclusion
(as given in the solutions on the internet). We begin by observing that since
a, b, c are in A.P.,

b−a= c−b (6)

Similarly, the fact that a2 , b2 , c2 are in H.P. translates into their reciprocals
1 1 1 1
being in A.P. and hence 2 − 2 = 2 − 2 , i.e.
b a c b
a2 − b2 b2 − c2
= (7)
a2 b2 b2 c2
Cancelling b2 , factoring the numerators and using (6), this finally reduces
to

(a − b)(c − a)(ac + 2b2 ) = 0 (8)

If either the first or the second factor vanishes that would mean a = b = c.
The vanishing of the third factor would mean that a, b and − 2c are in G.P.

43
Problem 12. Let a, b be arbitrary real numbers. Find the values of a for
which the equation x2 + (a − b)x + (1 − a − b) has unequal real roots for all b.

Analysis and Solution: This is a very straightforward problem based on

the criterion for a quadratic to have real roots. Taking the given expression
as a quadratic in x, the condition that it should have real and distinct roots
is that its discriminant be positive, i.e.

(a − b)2 > 4(1 − a − b) (1)

We want those values of a for which (1) will hold for all b. To do this we
rewrite the condition (1) as a quadratic expression in b being positive, i.e.

b2 + (4 − 2a)b + a2 − 4(1 − a) > 0 (2)

for all b. Now, for a fixed a, the L.H.S. is a quadratic expression in b with
leading coefficient 1 which is positive. The criterion for it to be positive for
all b is that it has no real roots. Or, taking discriminant,

(2 − a)2 < a2 − 4(1 − a) (3)

which simplifies to a > 1. So these are the values of a for which the given
equation has unequal real roots for all b.

Problem 13. Find the equation of the plane passing through the points
A(1, 2, 0), B(5, 0, 1) and C(4, 1, 1). Also determine the point Q such that the
line segment joining the points P (2, 1, 6) and Q is perpendicular to the plane
and is bisected by it.

Analysis and Solution: There are many ways of writing down the equation
of the plane passing through three given (non-collinear) points. They differ
more in their forms than in their mathematical contents. The simplest is to
apply the determinant criterion for the coplanarity of a general point (x, y, z)
and the three given points, say A, B, C. In the present problem this gives

x y z 1
 
 
 
 1 2 0 1 
=0 (1)
 
5 0 1 1
 
 
 
 4 1 1 1 

44
To cast this into the form ax + by + cz + d = 0 we have to evaluate four
3 × 3 determinants which is both time consuming and prone to errors. So
a short-cut is to take the general equation of a plane passing through one
of the given points and then to determine when it contains the other two
points. Taking this point as A(1, 2, 0), any plane passing through A has an
equation of the form

a(x − 1) + b(y − 2) + c(z − 0) = 0 (2)

where the constants a, b, c are to be determined. Superficially there are three

constants. But what matters is their relative proportions. So, to determine
them we need two conditions. These are given by the fact that the points
B(5, 0, 1) and C(4, 1, 1) lie on the plane. This translates into

4a − 2b + c = 0 (3)
and 3a − b + c = 0 (4)

Solving, a = b and c = −2a. So taking a (arbitrarily) as 1, (2) becomes

x + y − 2z = 3 (5)

Yet another method is to use vectors. Let P (x, y, z) be a typical point.

Then P will lie on the plane passing through A, B, C if and only if the
−→ −→ −→
vectors AP , AB and AC are coplanar, i.e. if and only if their box product is
−→
0. Writing these three vectors in terms of their components, AP = (x − 1)i +
−→ −→
(y − 2)j + zk, AB= 4i − 2j + k and AC= 3i − j + k and the box product as
a determinant, we get
 

 x − 1 y − 2 z 

 4 −2 1  = 0 (6)
3 −1 1 
 


which on expansion coincides with (5). Of course, we could also have gotten
(6) from (1) by substracting the second row from all others. (That is why,
the three methods are not conceptually very different.)
Now, for the second part of the problem, suppose P is the point (2, 1, 6)
and Q is the point such that the segment P Q is perpendicular to and bisected
by the plane. Then Q is the reflection of the point P in the plane. (Although
this interpretation does not simplify the calculations in the present problem,

45
such interpretations help in two ways. First, sometimes, the reflection of
a point in a plane can be found by inspection using some such features
as symmetry. Secondly, they also prepare you for questions involving such
reflections.)
To determine Q, we first observe that the vector i+ j−2k is perpendicular
−→ −→
to the plane represented by (5). (The vector AB × AC is also perpendicular
to the same plane. This knowledge helps in situations where we are not
interested in the equation of the plane itself but only in a vector perpendicular
−→
to it.) Hence P Q= λ(i + j − 2k) for some λ. This means Q = (2 + λ, 1 +
λ, 6 − 2λ). To determine λ, we use the condition that P Q is bisected by the
plane, or equivalently, that the midpoint, say M, of the segment P Q lies on
the plane. Clearly, M = (2 + λ2 , 1 + λ2 , 6 − λ). This satisfies (5) if and only if

λ λ
2+ + 1 + − 2(6 − λ) = 3 (7)
2 2
which gives λ = 4 and hence Q as (6, 5, −2).
Although the study of equations of straight lines and planes in three di-
mensional space is simple (and in many respects analogous to the study
of straight lines in a plane), it is often not formally covered at the Junior
College level. However, as the present problem shows, vectors can be used
conveniently to fill the consequent gaps.

Problem 14. Two tangents are drawn from the point P (6, 8) to the circle
x2 +y 2 = r 2 . Find r such that the area of the triangle formed by the tangents
and the chord of contact is maximum.

Analysis and Solution: This is a straightforward problem, consisting of

two parts, viz. (i) expressing the given area as a function of r and (ii)
maximising this function. However, as pointed out in Comment No. 2 of
Chapter 13, in maxima/minima problems a change of variable can often
simplify the calculations considerably. In the present problem, although we
are asked for which value of r the area is maximum, it is not mandatory to
express the area as a function of r. We are free to work with a variable of our
choice and after maximising the area we can find the corresponding value of
r.
In the present problem, it is in fact preferable to work in terms of the angle
θ which is half the angle which the chord of contact subtends at the point P

46
(see the figure). Let the tangents from P meet the circle at A and B. The
centre of the circle is at O, the origin. Clearly, OP = 10 and OA = OB = r.
So,
(6, 8) P
r = 10 sin θ (1) θ

which allows us to switch from r to A

θ and vice versa. From the figure
r
it is clear that
O
P A = P B = 10 cos θ (2) r
B
Denote the desired area, viz., the
area of the triangle P AB by f (θ).
Then we have
1
f (θ) = P A.P B sin 2θ = 100 cos3 θ sin θ (3)
2
where the variable θ ranges over the interval [0, π2 ]. The end-point 0 cor-
responds to the degenerate case where the circle reduces to a point while
the other endpoint π2 leads to the degenerate case where r = 10 and the
circle passes through the point P . f (θ) = 0 for both these values and so the
maximum must occur at an interior point of the interval [0, π2 ].
Now, differentiating (3) we get
f ′ (θ) = 100 cos4 θ − 300 cos2 θ sin2 θ = 100 cos2 θ(cos2 θ − 3 sin2 θ) (4)
Thus f ′ (θ) = 0 when either cos θ = 0 or when tan θ = √13 . The first case
gives the endpoint π2 at which f vanishes. The second case gives θ = π6 . So
this is the only critical point of f in the interval (0, π2 ). As we have already
shown that the maximum of f must occur at an interior point, no further
justification is needed to conclude that f attains its maximum on [0, π2 ] at
θ = π6 . Still, it is a good idea to also confirm this by observing from (4) that
f is increasing on [0, π6 ] and decreasing on [ π6 , π2 ].
To finsish the solution all we need to do is to find the corresponding value
of r. By (1) this value is 10 sin π6 = 5.
The solution given on the internet begins by expressing the area of the
!3/2
r2
triangle as a function of r. It comes out as 10r 1 − . Obviously,
100
this is a lot more complicated to work with than (3).

47
Problem 15. Normals with slopes m1 , m2 and m3 are drawn from a point
P , not on the axes, to the parabola y 2 = 4x. If the locus of P under the
condition m1 m2 = α is a part of the parabola, determine the value of α.

Analysis and Solution: This is a twisted locus problem. Normally in a

locus problem you are asked to find the locus of a point P which satisfies
some given condition. In the present problem, the condition is that the slopes
of two of the normals from P to the given parabola have a constant product
α. This is a different condition for different α’s and the locus too will change
as α varies. The problem asks you to determine that value of α for which
the locus will be a part of the parabola itself.
Thus, we have to begin by finding the locus of a point P = (h, k) which
satisfies the condition that the slopes of two of the normals from P to the
parabola y 2 = 4x have a constant product α. Let m be the slope of a normal
from (h, k) to the parabola y 2 = 4x. Then m must satisfy the cubic equation

k = mh − 2m − m3 (1)

(This is proved by assuming that the line y − k = m(x − k) is normal to

the parabola at some point (t2 , 2t) on it. But the equation to the normal at
the point (t2 , 2t) is y − 2t = −t(x − t2 ). Comparing these two equations of
the normal and eliminating t from it gives (1). Although the result is fairly
wel- known, it is not needed as frequently as the more standard results about
the parabola. This can lead to a controversy whether such a result can be
assumed without proof. In a problem like this where this result is only a
means and not the goal, it should be permissible to asume it without proof.)
(1) is a cubic in m and so in general it can have three roots, say m1 , m2 , m3 .
Their product is always −k. If, say m1 m2 = α, then m3 = − αk . So the
condition that the product of two roots of (1) be α is equivalent to saying
that − αk be a root of (1). (Note that since it is given that P does not lie
on the axes, none of the roots and hence the number α cannot be 0.) In in
other words,
hk 2k k 3
k=− + + 3 (2)
α α α
which upon simplification becomes k 2 = α2 h + α2 (α − 2). To get the locus
of P we replace h by x and k by y and get

y 2 = α2 x + α2 (α − 2) (3)

48
 In order that this locus be a part of the starting parabola y 2 = 4x, the
coefficients of the corresponding terms must be proportional. This implies
two equations
α2 = 4 (4)
and α2 (α − 2) = 0 (5)
The only solution of this system is α = 2, as we see by putting (4) into (5).
So this is the desired value of α.

Problem 16. Assuming that 2(1 − cos x) < x2 for all x 6= 0, or otherwise,
π
show that sin(tan x) ≥ x for all x ∈ [0, ].
4
Analysis and Solution: This is a problem where an inequality about func-
tions of a variable x has to be established by taking their derivatives. There
is a well-known inequality which says that
sin x ≤ x ≤ tan x (1)
for all x ∈ [0, π2 ). Since the sine function is also increasing on the interval
[0, π2 ], from the second inequality in (1) we get
sin(tan x) ≥ sin x (2)
So we would be through if we could show that sin x ≥ x. But in view of the
first inequality of (1) this attempt is doomed to fail. So, let us discard it
right away.
Let us, instead, define a function f (x) by
f (x) = sin(tan x) − x (3)
Then the problem amounts to showing that f (x) ≥ 0 for all x ∈ [0, π4 ].
Note that f (0) = 0. Therefore we would be through if we can show that for
all x ∈ [0, π4 ],
f ′ (x) ≥ 0 (4)
(A proof of this fact is given in Comment No. 11 of Chapter 13 using the
Rx
Mean Value Theorem. Or, we can write f (x) as f (x) − f (0) = f ′ (t) dt as
0
we did in the solution to Problem 8.)

49
 Now, by a direct calculation, f ′ (x) = cos(tan x) sec2 x − 1. Hence our
problem is reduced to showing that

cos(tan x) sec2 x ≥ 1 (5)

for all x ∈ [0, π4 ]. We are now in a position to use the hint given in the
problem. Put u = tan x. Then 0 ≤ u ≤ 1. Also, in terms of u, the L.H.S. of
(5) is simply (cos u)(1 + u2 ). The hint can be paraphrased to say that

u2
cos u > 1 − (6)
2
for all u 6= 0. (In fact, this is a more natural formulation of the hint. In this
form one can also see how to prove it. One method is to use the Taylor’s
theorem for the function g(u) = cos u. Of course, we are allowed to use the
hint without proof.)
u2
Because of this hint, we have (cos u)(1 + u2) ≥ (1 − )(1 + u2 ). Note
2
that we are also assuming that 0 ≤ u ≤ 1. So we shall be through if we can
u2
show that (1 − )(1 + u2 ) ≥ 1 for all u ∈ [0, 1]. This can be done by a
2
u2 u4
direct expansion of the product. It equals 1 + − . Now since u lies in
2 2
2
u u4
[0, 1], so does u2 , and therefore u2 ≥ u4 , whence − ≥ 0 which implies
2 2
2 4 2
u u u
1+ − ≥ 1, i.e. (1 − )(1 + u2 ) ≥ 1 as was to be shown.
2 2 2

Problem 17. Let In be the area of an n-sided regular polygon inscribed in

the unit circle and On be the area of an n-sided regular polygon circumscrib-
ing the same circle. Show that
 s 
2
On  2In

In = 1+ 1− 
2 n

Analysis and Solution: The notations In and On may seem strange. But
a moment’s thought will show that they correspond naturally to an ‘inner’ n-
gon and an ‘outer’ n-gon. An observation like this does not help you directly
in solving the problem. But it gives you a feeling that you are in tune with

50
the papersetter and that assurance can have a beneficial effect to boost your
confidence.
Now, actually coming to the solution, instead of hastily calculating In and
On , let us give it some thought. Obviously, for a fixed n, the regular n-gons in-
scribed in any two circles are similar to each other and the ratio of their areas
is the square of the ratio of the radii of the two circle. Moreover an escribed
regular n-gon can be thought of as an inscribed regular n-gon of a larger
circle. To see this, we draw a diagram
similar to that drawn in the solu- A3 B2
tion to Problem 14. A1 , A2 , A3 are
three consecutive vertices of the r R
regular n-gon inscribed in a cir- r A2
cle of radius r. Let the tangents
O π /n
at A1 and A2 meet at B1 while
r R B1
those at A2 and A3 meet at B2 .
Then it is clear that B1 and B2 A1
will be two consecutive vertices of
a regular polygon escribed about
this circle. This escribed n-gon
can also be thought of as a reg-
ular n-gon inscribed in a circle of
radius R where R = OB1 = OB2 .
To find In , we decompose the inscribed regular n-gon into n triangles each
of which is congruent to the triangle OA1 A2 . Since 6 A1 OA2 = 2π/n, we get
n 2 π π
In = r sin(2π/n) = nr 2 sin( ) cos( ) (1)
2 n n
In our problem, r is given to be 1. But for the moment we leave it as it is,
because by the same reasoning, we now get the area of the escribed regular
n-gon. The only change is that we replace r by R. Thus,
n 2 π π
On = R sin(2π/n) = nR2 sin( ) cos( ) (2)
2 n n
From the right-angled triangle OA2B1 in the figure we get R = r sec πn .
We now put r = 1. Then R = sec πn . Putting these values into (1) and (2)
respectively and dividing we have
2In 2 π
= π = 2 cos2 (3)
On sec2 n
n

51
 Also from the middle term of (1) (with r = 1), we have

2In 2π
= sin (4)
n n
from which it follows that
s
2
2In 2π π

1+ 1− = 1 + cos = 2 cos2 (5)
n n n
The assertion of the problem follows by combining (3) and (5).
In the internet solutions, the area of the triangle is calculated in an un-
necessarily complicateed way, by first finding the length of A1 A2 and then
that of the altitude through O. Our approach is simple because we already
know the sides OA1 , OA2 and the included angle. Similarly, in the internet
solution, the area of the triangle OB1 B2 is calculated separately. We have
bypassed this by observing that it is similar to the triangle OA1 A2 . Although
these things are relatively unimportant from a mathematical point of view,
in trigonometry we have a wealth of formulas and it pays to choose the one
that will give the answer effortlessly.

Problem 18. Let f : [0, 4] → IR be a differentiable function.

(a) Show that there exist a, b ∈ [0, 4] such that (f (4))2 −(f (0))2 = 8f ′ (a)f (b)

(b) Show that there exist α, β with 0 < α < β < 2 such that
Z 4
f (t)dt = 2(αf (α2) + βf (β 2))
0

Analysis and Solution: For (a), in absence of any other clue the only thing
that suggests itself is to factorise the L.H.S. as (f (4) + f (0))(f (4) − f (0)).
Now the second factor gives a clue. If we divide it by 4, which is the length
of the interval with end-points 0 and 4, we see that the ratio is amenable to
the Lagrange’s Mean Value Theorem. Thus we get that there exists some
a ∈ (0, 4) such that

f (4) − f (0)
= f ′ (a) (1)
4

52
Thus the problem is now reduced to showing that there exists some b ∈ [0, 4]
such that
f (0) + f (4)
= f (b) (2)
2
Here the trick is to note that the L.H.S. lies between the minimum, say m,
and the maximum, say M, of f (x) on the interval [0, 4]. (The existence of m
and M follows from continuity of f on the closed interval [0, 4].) This follows
by adding the two inequalities

m ≤ f (0) ≤ M
and m ≤ f (4) ≤ M

and dividing by 2. So, by the Intermediate Value Property of continuous

functions, we get some b ∈ [0, 4] for which (2) holds. As already proved, this
completes the solution of (a).
For part (b) of the problem, there is no obvious hint. But we can take
a cue form the expressions 2αf (α2 ) and 2βf (β 2) appearing in the R.H.S.
These look like the the values of the derivative of some function, say G(x),
at the points α and β respectively. So, we look for a function G(x) such that
G′ (x) = 2xf (x2 ). Moreover, since the derivative of G is realted to f , G ought
to be defined by an integral. Once this idea strikes, it is not difficult to come
up with such a function. Indeed one such function can be defined by
Z x2
G(x) = f (t)dt (3)
0

for then by the Fundamental Theorem of calculus (the second form), we have

G′ (x) = 2xf (x2 ) (4)

(Here we need the Fundamental Theorem coupled with chain rule as given
in Equation (19) of Chapter 17.)
R4
In terms of this new function G, the expression f (t) dt is nothing but
0
G(2) − G(0). If we divide this by 2 and apply Lagrange’s MVT, we get that
there exists some c ∈ (0, 2) such that
Z 4
f (t)dt = 2G′ (c) = 4cf (c2 ) (5)
0

53
 Unfortunately, this is not quite what we want. First, the R.H.S. of (5) has
the factor 4 instead of 2 which we would like to see there. Moreover, instead
of a single c, we want two numbers α and β. This suggests that we should
apply Lagrange MVT to the function G not on the single interval [0, 2] but
after dividing it into two halves. Doing so, we get that there exist α and β
in the intervals (0, 1) and (1, 2) respectively such that

G(1) − G(0) = G′ (α) = 2αf (α2) (6)

and G(2) − G(1) = G′ (β) = 2βf (β 2) (7)

R4
Adding (6) and (7) and noting that G(2) − G(0) is simply f (t) dt com-
0
pletes the solution of Part (b) of the problem, except that we have to ensure
that α < β. But this follows since we have α < 1 and 1 < β. The solution is
now complete to the hilt.
Considering that the Mean Value Theorems were introduced in the JEE
syllabus in 2003 after a long gap, the second part of the problem is a bit too
tricky. Having to split the interval [0, 2] into two halves makes it even more
so. It would have been reasonable to ask (5). Or, a subtle indication of the
need to split could have been given by putting the conditions 0 < α < 1 and
1 < β < 2, instead of the single condition 0 < α < β < 2 in the statement of
the problem.

Problem 19. Let ~u, ~v, w ~ be three non-coplanar unit vectors. Let α, β, γ be
the angles between the vectors ~u and ~v , ~v and w, ~ and w~ and ~u, respectively.
~
Let ~a, b, ~c be unit vectors along the bisectors of the angles α, β, γ respectively.
Show that
1 α β γ
[~a × ~b, ~b × ~c, ~c × ~a] = [~u, ~v , w]
~ 2 sec2 ( ) sec2 ( ) sec2 ( ),
16 2 2 2
where [~x, ~y , ~z] = ~x · (~y × ~z ).

Analysis and Solution: Obviously, the right place to begin is to get hold of
the vectors ~a, ~b, ~c, i.e. to express them in terms of the vectors ~u, ~v and w.
~ We
are given that ~a is a unit vector in the direction of the angle bisector of the
angle α between ~u and ~v . (In the old days, to keep perfect consistency with
cyclic symmetry, this angle would have been denoted by γ and the vector
~a by ~c. See the remark made in the solution to the problem in Comment

54
No. 13 of Chapter 7. Of course, mathematically that does not matter.)
Obviously, the vector ~a is sym-
v u+v
metrically inclined to ~u and ~v .
Moreover, ~u and ~v have the a
same lengths. So ~a = λ(~u + ~v )
for some (positive) scalar λ. To
determine it we use the fact α u
that ~a, ~u, ~v are unit vectors and
further ~u · ~v = cos α. So,

1 = |~a|2 = λ2 |~u + ~v |2
= λ2 (|~u|2 + 2~u · ~v + |~v|2 )
α
= λ2 (2 + 2 cos α) = 4λ2 cos2 ( )
2
1
which gives λ = 2
sec α2 and hence

1 α
~a = sec (~u + ~v ) (1)
2 2
By cyclic symmetry, we also have,

~b = 1 sec β (~v + w)
~ (2)
2 2
1 γ
and ~c = sec (w ~ + ~u) (3)
2 2
The next thing is to simplify the L.H.S. of the equality to be proved. By
a standard identity, for any three vectors ~a, ~b and ~c, we have

[~a × ~b, ~b × ~c, ~c × ~a] = [~a, ~b, ~c]2 (4)

(For a proof, see Exercise (21.2) (f). However, this is not such a frequently
used identity about the box product as its cyclic symmetry, for example. It
would have been better to give this as a hint or pose it as a separate part of
the problem.)
Because of (4), our problem is now reduced to proving that

1 α β γ
[~a, ~b, ~c] = [~u, ~v , w]
~ sec sec sec (5)
4 2 2 2

55
 ~ are unit vectors and ~a, ~b, ~c are given by (1) to (3). Using these
where ~u, ~v, w
to substitute for ~a, ~b, ~c, and by linearity of the box product in each argument,
we have
1 α β γ
[~a, ~b, ~c] = sec sec sec [~u + ~v , ~v + w,
~ w~ + ~u] (6)
8 2 2 2
The problem is now reduced to expanding the box product [~u +~v , ~v + w, ~ w+~
~ u].
Again, by its linearity in each argument it will split as the sum of eight box
products of the form [~x, ~y , ~z], where ~x is either ~u or ~v, ~y is either ~v or w
~ and
~z is either w ~ or ~u. Six of these eight box products will vanish because some
vector will be repeated in them. The only non-vanishing ones are [~u, ~v, w] ~
and [~v , w,
~ ~u] which are equal to each other by cyclic symmetry. Hence,

[~u + ~v , ~v + w,
~ w~ + ~u] = 2[~u + ~v + w]
~ (7)

(6) and (7) together imply (5) and as proved earlier, this completes the
solution.

Problem 20. A conical flask of height H has pointed bottom and circular
top of radius R. It is completely filled with a volatile liquid. The rate of
evaporation of the liquid is proportional to the surface area of the liquid in
contact with air, with the constant of proportionality K > 0. Neglecting
the thickness of the flask, find the time it takes for the liquid to evaporate
completely.

Analysis and Solution: There are many problems where liquid is drained
out from a container through a hole in the bottom. Here it is being lost
because of evaporation for a change. Still, since the rate of evaporation is a
dV
certain derivative (specifically, the derivative , where V is the volume of
dt
the liquid left and t is the time), we expect that this problem will involve
two things : (i) writing down a certain differential equation, and (ii) solving
it.

56
As with most problems where
the data is physical and geo- R
metric, drawing a diagram is
the correct start. The radius
of the top of the cone is given r
as R. Let r = r(t), h = h(t)
and V = V (t) denote, respec- V H
tively, the radius of the surface h
of the liquid left, its height and
its volume at time t (as mea-
sured from the start). Then
1
V = πhr 2 (1)
3
h r
By similarity of triangles, = . Putting this into (1) we get
H R
1 H
V = π r3 (2)
3 R
Differentiating this w.r.t. t gives
dV πH 2 dr
= r (3)
dt R dt
But we are given that at any time
dV
= −Kπr 2 (4)
dt
From (3) and (4) we get the desired differential equation for r, viz.
dr KR
=− (5)
dt H
As the R.H.S. is a constant, this differential equation is extremely easy to
solve. Integrating both the sides w.r.t. t we get
KR
r = r(t) = − t+c (6)
H
where c is a constant. At the start, i.e. at t = 0, the flask is full and hence
r(0) = R. This gives c = R. So, (6) becomes
KR
r = r(t) = R − t (7)
H
57
We are asked to find the time when the flask is empty, i.e. when r = 0. From
H
(7), this happens when t = K .
We have got the answer systematically to illustrate the procedure. How-
ever, in the present case, the utter simplicity of the differential equation (5)
makes it possible to arrive at the answer by sheer common sense. For, (5)
can be interpreted to say that the radius r is decreasing at a constant rate,
viz. KR
H
. Hence to go down from the initial value R to the final value 0, the
H
time taken will be R/(KR/H) = K .
Because of this alternate solution, the problem is essentially a problem
about derivatives, specifically about related rates. The aspect of solving a
differential equation is hardly tested in this problem.
It is interesting to see the physical significance of the result. The constant
K depends upon the liquid. So the answer we got says that the time taken
for the given liquid to evaporate depends only on the height and not on the
shape of the cone. So, if you want to store it for long in a conical flask, make
the depth of the flask as big as you can by taking its top as small as you can.

58
 CONCLUDING REMARKS

Note: These remarks are meant to take an overall view of the Screening and
the Main Papers as opposed to a problem-wise one. The opinions expressed
are strictly personal opinions of the author.

The screening paper is quite reasonable both in terms of the types of

the questions and the work needed to solve them. Hoever, as mentioned in
the comments on several questions (for example, Q. 1 and Q. 18), a little
care could have been exercised to preclude unintended short-cuts. There is
some duplication of ideas in Q. 21 and Q. 23 since they both involve the
intercepts made on the coordinate axes by a tangent to an ellipse. One of
these problems could have been replaced by a problem in Number Theory,
which has not at all been catered to. Except for this omission, the Screening
Paper covers all the syllabus as completely as can be expected from multiple
choice questions.
It is difficult to give an equally clean chit to the Main Paper. Here too,
the omission of Number Theory is shocking. It is true that many problems in
number theory can be tackled by induction and ‘induction’ is no longer a spe-
cific topic listed in the syllabus. But that should only mean that questions
designed only to test the knowledge of induction cannot be asked. It cer-
tainly should not mean that induction is outlawed as a tool, because in that
case many interesting problems (including many trigonometric identities, bi-
nomial identities, reduction formulas for integrals and certain problems in
probibility) will have to be discarded.
It can perhaps be argued that in order to make room for new topics like
matrices and Mean Value Theorems introduced into the JEE syllabus, the
axe has to fall on some old topics. Even then, it is difficult to justify two
problems (1 and 2) based on inequalities about complex numbers. Surely,
one of them could have been replaced by a Number Theory problem. There
are also other areas (e.g. differential equations, probability) where more than

59
one question are asked each. There are two maxima/minima problems, viz.
Problem 4 and Problem 14. The latter could have been dispensed with since
the part involving the calculation of the area in it is essentially subsumed
in Problem 17 which is a good problem. The probability problems are also
good. In fact, the first one (Problem 5) is innovative in that the trials are
not independent as is usually the case. This is probably the best problem in
the entire Main Paper.
Having to solve 20 problems in two hours seems unreasonable, especially
so considering that in two of these problems (Problem 3 and Problem 8)
there was some inconsistency in the data of the problem. Needless to say,
such mistakes are very unsettling to the intelligent students. They spend
a lot of their precious time to check and recheck if they have made some
mistakes. It is high time that a uniform policy is laid about such erroneous
questions and further that such a policy is made known to the candidates
well in advance.
Considering the overall nature of the Main Paper, three hours would have
been the right time for solving it. In fact, it would be far better in future to
have two mathematics papers of 2 to 2.5 hours each, catering separately to
pre-calcus mathematics and calculus-based mathematics. That will ensure
that no part of the syllabus is left untouched. A similar case be made for the
other two subjects, viz. Physics and Chemistry. In that case the examination
will have to spread over three days instead of finishing it all in a day as at
present. But what is the harm? Considering that only a limited number
of candidates appear for the Main Examination which is conducted at only
a few centres over the country, the task of organising it and evaluating six
papers per candidate is not such a time consuming task. In fact, this will
give a more relaxed and valid testing of a candidate’s abilities. The present
practice of testing a candidate in three separate subjects in a single day for
six hours is more a test of stamina than of other more important qualities.
A comment is also perhaps in order about the solutions displayed on the
internet. The transparency revealed in adopting this practice is highly com-
mendable. However, to make it even more so, it would be desirable if alter-
nate solutions are also given wherever possible. Although it is mentioned at
the start of the solutions that alternate solutions, if any, have also been con-
sidered, giving only one solution on the internet often gives an unwarranted
impression that it is the standard or the expected solution, even though
sometimes better solutions are possible, as we have pointed out in the case
of Problems numbered 3, 4, 10, 11, 13, 14 and 17.

60
 The year 2003 is the very first year when the organisers of the Joint En-
trance Examination have taken the step of making the questions as well as
the answers available on the internet. This is a highly commendable gesture
which will vindicate the public faith in the JEE. The shorcomings pointed
above will hopefully be eliminated in the years to come.

61